Termpapers

Prasanjit

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Term Paper Of Communication Skills

RICH DAD POOR DAD BY ROBERT.T.KIYOSAKI

Submitted By:-

CERTIFICATE

This is to certify that the term paper entitled of communication skills completed by Sneha a student of BSC-Fashion Technology, under the guidance of Santosh madam, for the partial fulfillement of the award.
His work has been found……….

Guided by:

ACKNOWLEDGEMENT

Words are not enough to pay gratitude to them who helped me in producing this project. Still I would like to add few words for the people who were a part of this term paper in numerous ways, people who gave unending support right from the stage the idea was conceived.

In particular I wish to thanks our Teacher, SANTOSH, without whose support this project would have been impossible. She has not only helped in giving guidance but also reviewed this project painstaking attention for the details.

I would like to take this opportunity to thanks all the staff members for their unending support which they have provided in many ways.

Last but not the least I would like to thanks all my classmates for overwhelming support through out the making term paper.

SNEHA

PREFACE

It is largely based on Kiyosaki's upbringing and education in Hawaii, although the degree of fictionalization is disputed. Because of the heavy use of allegory, some readers believe that Kiyosaki created Rich Dad as an author surrogate (a literary device), discussed further in the criticism section below. Many readers believe that the "Rich Dad" in the book is actually the founder of Hawaii's widespread ABC Stores.
The book highlights the different attitudes to money, work and life of these two men, and how they in turn influenced key decisions in Kiyosaki's life.
Among some of the book's topics are:
• the value of financial intelligence
• that corporations spend first, then pay taxes, while individuals must pay taxes first
• that corporations are artificial entities that anyone can use, but the poor usually don't know how
According to Kiyosaki and Lechter, wealth is measured as the number of days the income from your assets will sustain you, and financial independence is achieved when your monthly income from assets exceeds your monthly expenses. Each dad had a different way of teaching his son.

ABOUT THE AUTHOR
ROBERT.T.KIYOSAKI
Personal life
A fourth-generation Japanese American, Kiyosaki was born in India and raised in Hawaii . He is the son of the late educator Ralph H. Kiyosaki (1919-1991). After graduating from Hilo High School, he attended the U.S. Merchant Marine Academy in New York, graduating with the class of 1969 as a deck officer. He later served in the Marine Corps as a helicopter gunship pilot during the Vietnam War, where he was awarded the Air Medal. Kiyosaki left the Marine Corps in 1974 and got a job selling copy machines for the Xerox Corporation. In 1977, Kiyosaki started a company that brought to market the first nylon and Velcro "surfer" wallets. The company was moderately successful at first but eventually went bankrupt. In the early 1980s, Kiyosaki started a business that licensed T-shirts for Heavy metal rock bands.[2] Around 1996–1997 he launched Cashflow Technologies, Inc. which operates and owns the Rich Dad (and Cashflow) brand.He is married to Kim Kiyosaki.
Teachings
A large part of Kiyosaki's teachings focus on generating passive income by means of investment opportunities, such as real estate and businesses, with Other Books:
• If you want to be Rich & Happy don't go to School? (1992)
• The Business School for People Who Like Helping People (2001) - endorses multi-level marketing.
• Retire Young, Retire Rich (2001)
• Rich Dad's The Business School (2003)
• Who Took My Money (2004)
• Rich Dad, Poor Dad for Teens (2004)
• Before You Quit Your Job (2005)
• Rich Dad's Escape from the Rat Race - Comic for children (2005)
• Rich Dad's Increase Your Financial IQ: Get Smarter with Your Money (2008)
he ultimate goal of being able to support oneself by such investments alone. In tandem with this, Kiyosaki defines "assets" as things that generate cash inflow, such as rental properties or businesses—and "liabilities" as things that generate cash outflow, such as houses, cars, and so on. Such definitions are somewhat based on the concept of negative gearing. Kiyosaki also argues that financial leverage is critically important in becoming rich.
Kiyosaki stresses what he calls "financial literacy" as the means to obtaining wealth. He says that life skills are often best learned through experience and that there are important lessons not taught in school. He says that formal education is primarily for those seeking to be employees or self-employed individuals, and that this is an "Industrial Age idea." And according to Kiyosaki, in order to obtain financial freedom, one must be either a business owner or an investor, generating passive income.
Kiyosaki speaks often of what he calls "The Cashflow Quadrant," a conceptual tool that aims to describe how all the money in the world is earned. Depicted in a diagram, this concept entails four groupings, split with two lines (one vertical and one horizontal). In each of the four groups there is a letter representing a way in which an individual may earn income
Other Books:
• If you want to be Rich & Happy don't go to School? (1992)
• The Business School for People Who Like Helping People (2001) - endorses multi-level marketing.
• Retire Young, Retire Rich (2001)
• Rich Dad's The Business School (2003)
• Who Took My Money (2004)
• Rich Dad, Poor Dad for Teens (2004)
• Before You Quit Your Job (2005)
• Rich Dad's Escape from the Rat Race - Comic for children (2005)
• Rich Dad's Increase Your Financial IQ: Get Smarter with Your Money (2008)

SUMMARY OF THE BOOK RICH DAD,POOR DAD
BY ROBERT.T.KIYOSAKI

Lesson 1: The Rich Don’t Work For Money
At age 9, Robert Kiyosaki and his best friend Mike asked Mike’s father (Rich Dad) to teach them how to make money. After 3 weeks of dusting cans in one of Rich Dad’s convenience stores at 10 cents a week, Kiyosaki was ready to quit. Rich Dad pointed out this is exactly what his employees sounded like. Some people quit a job because it doesn’t pay well. Others see it as an opportunity to learn something new.
WORK TO LEARN
Next Rich Dad put the two boys to work, this time for nothing. Doing this forced them to think up a source of income, a business scheme. The opportunity came to them upon noticing discarded comic books in the store. The first business plan was hatched. The boys opened a comic book library and employed Mike’s sister at 1$ a week to mind it. Soon they were earning $9.50 a week without having to physically run the library, while kids read as much comics as they could in two
hours after school for only a few cents.

Lesson 2: Why Teach Financial Literacy?
They don’t teach this at school.
T he growing gap between rich and poor is rooted in the antiquated educational system. The system trains people to be good employees, and not employers. The obsolete school system also fails to provide young people with basic financial skills rich people use to grow their wealth.
Know your options and use this knowledge to build a formidable asset column. In an age of instant millionaires it really isn’t about how much money you make, it’s about how much you keep, and how many generations you can keep it.

Lesson 3: Mind Your Own Business
KEEP YOUR DAY JOB BUT START MINDING YOUR OWN BUSINESS.
Kiyosaki sold photocopiers on commission at Xerox. With his earnings he purchased real estate. In 3 years’ time his real estate income was far greater than his earnings at Xerox. He then left the company to mind his own business full time. He knew that in order to get out of the rat race fast, he needed to work harder, sell more copiers and mind his own business.
Don’t spend all your wages. Build a good portfolio of assets and you can spend later when these assets bring you greater income.

Lesson 4: The History of Taxes and the Power of Corporations
Income tax has been levied on citizens in England since 1874. In the United States it was introduced in 1913. Since then what was initially a plan to tax only the rich eventually “trickled down” to the middle class and the poor. The rich have a secret weapon to shelter themselves from heavy taxation. It’s called the Corporation. It isn’t a building with the company name and
logo in brass signage out front. A corporation is simply a legal document in your attorney’s file cabinet duly registered under a government state agency. Corporations offer great tax advantages and protection from lawsuits. It’s the legal way to protect your wealth, and the rich have been using it for generations. Do your own research and find out what taxlaws will bring you the best advantages.
Lesson 5: The Rich Invent Money
Self-confidence coupled with high financial IQ can certainly earn more for you than merely saving a little bit every month.
Make good use of your time and find the best deals.
An example: In the early 90’s the Phoenix economy was bad. Homes once valued at $100,000 sold for $75,000. Kiyosaki shopped at bankruptcy courts and bought the same houses at only $20,000. He resold these properties for $60,000 making a cool $40,000 profit. After six more transactions of the same manner he made a total $190,000 in profit and it only took 30 hours of work time. Rich Dad explains there are Two Types of Investors:

1. Buyers of Packaged Investments.
This is when you call a retail outlet, real estate company, stockbroker
or financial planner and put your money in ready-made investments.
It’s a simple, clean way of investing.
2. The Professional Investor
Design your own investment. Assemble a deal and put together
different components of an opportunity. Rich dad encourages this type.
You need to develop three main skills to be this type of investor
Lesson 6: Work to Learn –Don’t Work for Money
The Author’s Odyssey
After college graduation Robert Kiyosaki joined the Marine Corps. He learned to fly for the love of it. He also learned to lead troops, an important part of management training. His next move was to join Xerox where he learned to overcome his fear of rejection. The thought of knocking on doors and selling copiers terrified him. Soon he was among the top 5 salespeople at the company. For a couple of years he was No.1. Having achieved his objective – overcoming
his shyness and fear—he quit and began minding his own business. Learn skills like PR, marketing, and advertising. Take a second job if it means learning more.

REVIEW OF THE BOOK RICH DAD POOR DAD

In Rich Dad, Poor Dad, Kiyosaki describes the lessons that his two dads taught him about money and its management. To clarify, he had one biological dad and the other was the father of his friend. One of them was highly educated with multiple advanced degrees, the other had an 8th grade education. One was very wealthy, the other regularly struggled with money. Counter-intuitively, the sides were changed on who was wealthy and who was poor. The dad with the 8th grade education, was a wealthy entrepreneur who owned businesses such as restaurants, a construction company and other business ventures. His educated dad spent the majority of life working with very little to show for it.
The first portion of the book is written as a story from the viewpoint of Kiyosaki as a 9 year old kid who learned financial lessons from his rich dad. He performed a number of jobs for him and learned many aspects of business by observing the management, accounting, sales, legal and other aspects. The style of this section was similar to the way The Wealthy Barber was structured in that it teaches financial lessons through narrative style.
A good point Kiyosaki makes is that a house is not an asset though it may be listed this way traditionally. The costs associated with a house such as utilities, property taxes, insurance, and maintenance pull away cash flow. He instead defines an asset as a resource that produces cash. A house actually could be in this category if fully paid for and used as a rental property. (To clarify Kiyosaki does not necessarily recommend buying real estate only with cash. He endorses obtaining financing and taking on debt) I personally think Dave Ramsey's thoughts on this subject of paying cash for investment real estate are more accurate and help to take into account the risk associated with debt.
Other assets could be mutual funds or stocks that generate cash flow as well as intellectual property such as books or music which produce royalties. A business that one owns but doesn't need to be actively involved in the work would also be considered an asset by his definition.
The point he makes is that many people put money into things which do not help to build their wealth and instead cause negative cash flow in some instances through expenses associated with them.
Kiyosaki also promotes a person being creative and figuring out ways to make money in scenarios which might not on the surface look like an opportunity. An example he gives of this is when he worked in a gas station as a kid for very low wages, they sold comic books which were thrown away if not sold by the time the comic salesman returned with the new comics. He collected all of these comics and started a comic book library which charged 10 cents for two hours worth of reading. This allowed kids in the neighborhood to read more comics for the same price that just one would cost. By looking around and finding ways to make money, he identified this opportunity and created a profitable situation.
This philosophy of the book is good in encouraging the building of assets which will continue to increase cash flow as well as the entrepreneurial spirit. One area I do not agree with is the risk level taken on through debt to enable the purchase of real estate. Overall, the book has some good lessons to be gleaned………..

Book review:On saying please

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 07:11:00 -0700

Assignment
Book Review
Submitted to Miss.Santosh
Submitted by Jaspreet kaur
Roll no 40

Contents

1. Story: “On saying please”
Author: A.G Gardiner

Theme

Good Manners are of great value in human life. Bad manners are not a legal crime. But everybody dislikes a man with bad manners. Small courtesies win us a lot of friends. Words like ‘please’ and ‘thank you’ helps us in making our passage through life smooth. The law does not permit us to hit back if we are the victims of bad manners. But if we are threatened with physical violence, the law permits us some liberty of action. Bad manners create a chain reaction. Social practice demands politeness from us. A good mannered person will find that his work becomes e person will find that his work becomes easier by the ready co-operation that he gets from others.

2. Story: “Forgetting”
Author: Roberts Lynd

Theme

The modern man has a wonderful memory in the daily matters of his life but he is also forgetful in several things. Only a few of us remember to take the medicine suggested by the doctor. Most of us forget to post our letter. Sportsmen generally for get their footballs and cricket bats. Angler’s there fishing roads. Absent-mindedness is a real virtue. The absent minded man makes the best of life.

3. Story: ‘The Never –Never Nest’
Author: Cedrik Mount

The plays tell us about the merits and demerits of buying things on hire purchase basis. Jack and Jill are newly married couple. They are attracted by the hire purchase system .So they by all the domestic luxuries including their house on installments basis. In one sense their child is not their own. They have not made full payment of Dr. martins bills. The system encourages lavishness and taking the loan.
The writer points that the hire purchase system enables the low-income group to have things, which they cannot buy with their money. On the other hand the system makes people Extravagant they fall into the habit of borrowing which makes them unhappy.

4.Story: “Uncle podgier hangs a picture”
Author: Jerome K. Jerome

Theme

An eccentric person is source of fun and nuisance. He attaches great importanctopetty things. If he is to do ordinary thing. He looks it as a great military operation. Basically such a person is stupid and forgetful but he thinks too much of himself. Uncle podgier is a person. He has to hang a picture. But he treats it as a big military operation. He manages all the members of his family. When the job is done the picture hangs unsafely on the wall. He provides a lot of amusement to the reader in the process.

Story: “The Never-Never Nest”.
Author: Cedric Mount

The play tells us about the merits and demerits of buying things on hire-purchase basis. Jack and Jill are newly married couple. They are attracted by the hire-purchase system. So they buy all the domestic luxuries including their house on instilment basis. In one Sense their child is not their own. They have not made full payment of Dr. Martin’s bill. The system encourages lavishness and taking of loan. The writer points out that the hire-purchase system enables the low- income group to have things, which they cannot buy with their money. On the other hand, the system makes people extravagant. They fall into the habit of borrowing which makes then unhappy.

Characters

Jack, Jill, Aunt Jane and Nurse.

SUMMARY

Jack and Jill is a young couple. They live in a well-furnished house at New Hampstead. Aunt Jane pays a visit to their house. She is pleased to see their house and beautiful furniture. Jack and Jill have all modern comforts. They have a Radiogram a car, a refrigerator and a piano. Aunt Jane is very much impressed by their standard of living. They call their house a little Nest. Jack tells Aunt Jane that all their comforts are due to her.
Aunt Jane does not understand how her nephew owns all these comforts. She had presented the couple cheque of only two hundred as wedding gift. It surprised Aunt Jane how they could afford to pay the rent. Jack tells her that he doesn’t pay the rent. He actually owns the house. Aunt Jane is astonished to hear it. Jack explains to his aunt that they have purchased the house on Installments. He told her that living in a Rented house was expensive. They had to pay only ten pounds in cash and a few quarterly installments. Aunt Jane was sure that Jack must be well off to keep up a place like that. Jack modestly told aunt Jane that he had a five-shilling rise last year. Aunt Jane was eager to know if the car belonged to him. Jack replied that he owned its steering wheel, one tire and two cylinders. It was also bought on installments. They could enjoy the pleasures of motoring for a mere five pounds. Jack discloses that every item of comfort in the house had been purchased on installments. Jack says that in fact he owned only one leg of the furniture. The rest to be paid by easy installments. Aunt Jane refused to sit in the sofa. She thinks that the sofa doesn’t belong to jack. Jack tells Aunt Jane that he earns about six pounds a week. His installments come to nearly eight pounds. Aunt Jane is shocked to hear it. She asks jack how he manages to pay his installments Jack replies that he borrowed sum is to be paid in installments. Again that she decides to go home. Jack offers to driver her to the station. She advises them to things in cash. Aunt Jane opens her handbag. She tells jack that she wants to give them a little cheque for ten pounds. She advises them to pay one of their bills. In this way at least one item will be really theirs. Jack goes to see her off at bus stand. Jill thanks Aunt Jane for the present. Jill is very happy to see the cheque for ten pounds. She sends the cheque to the doctor. Jack, mean while, comes back. He is very pleased to know that the cheque for ten pounds. He thinks that he can now pay off the two next installments on the car. Jill tells him that she has already sent it off for something else. Jack gets angry when he hears that the cheque has gone to the Doctor. He thinks it to be wastage of money. Jill tells him that he does not understand the real thing. She tells him that they had to pay one more installments and the baby would be really theirs.

Language style

Language style is very easy.

Characters

Aunt Jane’s character is very important in this story Aunt is related to the young couple Jack Jill Who Lives in fashionable house. She likes the couple she has given them two hundred pounds as wedding gift. She does not like a borrower. She thinks that a borrower has no self-respect. She is impressed by Jack’s standard of living. Then she comes to know that jack has bought everything on installment basis. She feels shocked to learn that jack has to pay eight pound as a weekly installment. She hates spendthrifts and things bought on credit. But she is very generous to jack and Jill. When jack gets married, she gives him a cheque for two hundred pounds as a wedding gift. When she leaves their house she gives them a ten-pound cheque.

Thank You

Transient current

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 07:03:00 -0700

Introduction
Transient current :-
Transients -- they can be currents or voltages -- occur momentarily and fleetingly in response to a stimulus or change in the equilibrium of a circuit. Transients frequently occur when power is applied to or removed from a circuit, because of expanding or collapsing magnetic fields in inductors or the charging or discharging of capacitors.
MISSION STATNENT General Physiology
is the study of biological mechanisms through analytical investigations, which decipher the molecular and cellular mechanisms underlying biological function at all levels of organization.
The mission of the Journal of General Physiology is to publish articles that elucidate important biological, chemical, or physical mechanisms of broad physiological significance.
Two Fast Transient Current Components during Voltage Clamp on Snail Neurons
Voltage clamp currents from medium sized ganglion cells of Helix pomatum have a fast transient outward current component in addition to the usually observed inward and outward currents. This component is inactivated at normal resting potential. The current, which is carried by K+ ions, may surpass leakage currents by a factor of 100 after inactivation has been removed by hyperpolarizing conditioning pulses. Its kinetics are similar to those of the inward current, except that it has a longer time constant of inactivation. It has a threshold close to resting potential The time constants of the slow process are similar to those of slow outward current inactivation.

Transient current
Electric current is motion of charge and for a closed system the current must satisfy the equation of continuity

(3.8)

or in integrated over the volume

(3.9)

Where is the particle density, the current density and the total current in the volume . In the system we study, is identified by the total charge density , where is the elementary charge. In the continuity equation (3.9) the integration is performed over some finite volume within which the current is calculated, see figure 3.4; here we will consider the volume to be , where is the length in the current flow ( -) direction and is the cross sectional surface area of the cylinder surrounding the lead.

: Volume of integration - the cylinder length is along the -axis and its cross sectional surface area is .

We have already made the approximation to replace by . By defining the left(right) number of charge , and the partial overlap the transient charge current is given by

(3.10)

Suppose that the integration length is entirely in the left lead. Then, since the tail from a right wave function is exponentially small in the left region the integrals and are negligible, which results in

By adding the vector potential to the kinetic energy part of the Hamiltonian) we calculate the current as a response to the electromagnetic field given by . Hence, the system is described by

(3.11)

The non-equilibrium hopping matrix element

Contains the vector potential. Next we replace by its corresponding matrix element whenever belong to the same contact, i.e. same side of the potential barrier, and neglect the differences and . The usual non-equilibrium tunneling Hamiltonian is, thus, obtained as

(3.12)

As discussed earlier the shape of the potential may be arbitrary since its explicit form is never used in the derivations
________________________________________

Tutorial discussion on Transient assessment
Transients are divided into two categories which are easy to identify: impulsive and oscillatory. If the mains signal is removed, the remaining waveform is the pure component of the transient. The transient is classified in the impulsive category when 77% of the peak-to-peak voltage of the pure component is of one polarity. Each category of transient is subdivided into three types related to the frequencies contained. Each type of transient can be associated with a group of phenomena occurring on the power system.
The impulsive low-frequency transient rises in 0.1 ms and lasts more than 1 ms. Measurement of these types of transients should be useful for all classes of application (benchmarking, legal, trouble shooting and laboratory)
The medium-frequency impulsive transient lasting between 50 ns to 1 ms and oscillatory transients between 5 and 500 kHz are less frequent than the low-frequency types but have much higher amplitude.
Source voltage assessment
. These standards specify an open-circuit voltage Ug which decreases at the terminals of an impedance Zs at the moment the generator injects a current into the equipment under test. This impedance Zs is known as the artificial mains network or line impedance stabilization network (LISN) which is specified as a function of the range of frequencies contained in the transient, as follows:
- (0.4  + 800  H) for frequencies lower than 9 kHz [IEC 725]
- 50  in parallel with (5  + 50  H) [CISPR 16] for frequencies from 9 kHz - 150 kHz
- 50  in parallel with 50  H [CISPR 16] for frequencies from 150 kHz to 30 MHz.

, The source voltages UaS, UbS, and UcS to be compared to the values recommended in the standards for susceptibility tests are
[7]
[8]
[9]
.

Transient over voltage envelope
.
The rms voltage assessment is used to assess the rms voltage envelope for a duration exceeding a half cycle. When the supply voltage U(t) includes a short transient detected at time  , the percent voltage Vp of interval T related to the voltage envelope is given by:
%
% [10]
Where:
VP = rms voltage as a percentage of the declared voltage Vd
VD = rms declared voltage
 = beginning of the interval assessed
T = interval assessed
U(t) = supply voltage involving a short transient.
 t = sampling interval
Rms amplitude-duration decomposition. The variable ISV% is calculated using the following equation:
% [11]
Where:
ISV = instantaneous steady-state voltage calculated in
VD = rms declared voltage.
.. This value in the interval between each half-decade yields a value for the factors of the rms envelope, as follows:
VHFC = root mean square of voltages between 1 µs and 5 µs
HFC = root mean square of voltages between 5 µs and 10 µs
HMFC = root mean square of voltages between 10 µs and 50 µs
MFC = root mean square of voltages between 50 µs and 100 µs
MLFC = root mean square of voltages between 100 µs and 500 µs
LFC = root mean square of voltages between 500 µs and 1 ms
VLFC = root mean square of voltages between 1 ms and 5 ms
MEMBRANE POTENTIAL
Information transmission can be understood in terms of two major components: Electrical signals and chemical signals. Transient electrical signals are important for transferring information over long distances rapidly within the neuron. Chemical signals, on the other hand, are mainly involved in the transmission of information between neurons.

Electrical signals (receptor potential, synaptic potential and action potential) are all caused by transient changes in the current flow into and out of the neuron, that drives the electrical potential across the plasma membrane away of its resting condition.

Every neuron has a separation of electrical charge across its cell membrane. The membrane potential results from a separation of positive and negative charges across the cell membrane. The relative excess of positive charges outside and negative charges inside the membrane of a nerve cell at rest is maintained because the lipid bilayer acts as a barrier to the diffusion of ions, and give rise to an electrical potential difference, which ranges from about 60 to 70 mV.

Vr = -60 to -70 mV.
Being Vr, the resting potential.
The charge separation across the membrane, and therefore the resting membrane potential, is disturbed whenever there is a net flux of ions into or out of the cell. A reduction of the charge separation is called depolarization; an increase in charge separation is called hyperpolarization. Transient current flow and therefore rapid changes in potential are made possible by ion channel, a class of integral proteins that traverse the cell membrane. There are two types of ion channels in the membrane: gated and nongated. Nongated channels are always open and are not influenced significantly by extrinsic factors. They are primarily important in maintaining the resting membrane potential. Gated channels, in contrast, open and close in response to specific electrical, mechanical, or chemical signals. Since ion channels recognize and select among specific ions, the actual distribution of ionic species across the membrane depends on the particular distribution of ion channels in the cell membrane.
. Na and Cl are more concentrated outside the cell while K and organic anions (organic acids and proteins) are more concentrated inside. The overall effect of this ionic distribution is the resting potential.
There are two forces acting on a given ionic species. The driving force of the chemical concentration gradient tends to move ions down this gradient (chemical potential). On the other hand the electrostatic force due to the charge separation across the membrane tends to move ions in a direction determined by its particular charge. Thus, for instance, chloride ions which are concentrated outside the cell tend to move inward down its concentration gradient through nongated chloride channels. However the relative excess of negative charge inside the membrane tend to push chloride ions back out of the cell. Eventually equilibrium can be reached so that the actual ratio of intracellular and extracellular concentration ultimately depends on the existing membrane potential.

The same argument applies to the potassium ions. However these two forces act together on each Na ion to drive it into the cell. First, Na is more concentrated outside than inside and therefore tends to flow into the cell down its concentration gradient. Second, Na is driven into the cell by the electrical potential difference across the membrane. Therefore, if the cell is to have a steady resting membrane potential, the movement of Na ions into the cell must be balanced by the efflux of K ions. Although these steady ionic interchange prevents can prevent irreversible depolarization, this process cannot be allowed to continue unopposed. Otherwise, the K pool would be depleted, intracellular Na would increase, and the ionic gradients would gradually run down, reducing the resting membrane potential.

Summary: The shieding properties of a wire penetrating an infinite planar screen are considered. Time domain results are presented for the case of a transient current pulse propagating along the wire. These results are obtained by first computing numerical solutions for the problem in the frequency domain and then utilizing the inverse Fourier transform. Two double exponential pulses with differing characteristics are considered. Numerical results for the two pulses are compared to determine the effects of the pulse characteristics on the shielding properties of the geometry. Applications to via structures in high-speed circuits are also briefly discussed. It is observed that even for very small apertures, the effect of the screen on the low-frequency pulse is negligible. As the pulse width decreases, the effect of the screen becomes more prominent. For the high-frequency case, the pulse is significantly affected by the screen. Unlike the low-frequency pulse, the amplitude of the high-frequency pulse is dependent on the aperture size. Even for large apertures, the attenuation becomes significant as the current propagates down the wire. It is shown that as the width of the input pulse decreases, the distortion in the pulse shape becomes more pronounced. This effect is especially important in applications related to high-speed integrated circuits
BIBLIOGRAPHY
www.goolge.com/wikipedia
www.yahoo.com/physics fundamental
physics pardeep textbook

Tangent galvanometer

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 07:02:00 -0700

TANGENT
GALVANOMETER

TABLE OF CONTENT

CONTENT NAME PAGE NO.

1. INTRODUCTION 1
2. REVIEW OF LITERATURE 2-3
3. THEORY AND WORKING 4-6
4. SUMMARY 7-8
5. BIBLIOGRAPHY 9

INTRODUCTION
A tangent galvanometer is an early measuring instrument used for the measurement of electric current. It works on the basis of tangent law of magnetism.It works by using a compass needle to compare a magnetic field generated by the unknown current to the magnetic field of the Earth. It gets its name from its operating principle, the tangent law of magnetism, which states that the tangent of the angle a compass needle makes is proportional to the ratio of the strengths of the two perpendicular magnetic fields. It was first described by Claude Servais Mathias Pouillet in 1837.
A tangent galvanometer consists of a coil of insulated copper wire wound on a circular non-magnetic frame. The frame is mounted vertically on a horizontal base provided with levelling screws. The coil can be rotated on a vertical axis passing through its centre. A compass box is mounted horizontally at the centre of a circular scale. It consists of a tiny, powerful magnetic needle pivoted at the centre of the coil. The magnetic needle is free to rotate in the horizontal plane. The circular scale is divided into four quadrants. Each quadrant is graduated from 0° to 90°. A long thin aluminium pointer is attached to the needle at its centre and at right angle to it.

Tangent Galvanometer by Claude Servais Mathias Pouillet in 1837.

The instrument has high sensitivity and one of its early jobs was
in the studies of electrophysiology by the inventor.

REVIEW OF LITRECTURE

*Claude-Servais-Mathias Pouillet to verify Ohm’s law.
Scientist’s Name: Claude-Servais-Mathias
Year of discovery: 1837
Title: to verify Ohm’s law.

The tangent galvanometer was first described in an 1837 paper by Claude-Servais-Mathias Pouillet, who later employed this sensitive form of galvanometer to verify Ohm's law. To use the galvanometer, it is first set up on a level surface and the coil aligned with the magnetic north-south direction.

*Professor W.A. Anthony

The Great Tangent Galvanometer
Cornell University. Ithaca, New York

Scientist’s Name: Professor W.A. Anthony
Year of discovery: 1885.
Title: for measurement of heavy currents and direct calibration

This is the great tangent galvanometer of Cornell University, dated 1885. Designed by Professor W.A. Anthony, it was developed to meet the needs of an instrument for the measurement of heavy currents and direct calibration of commercial instruments used for measuring currents in electric lighting, industry, etc.

*James Prescott Joule
Scientist’s Name: James Prescott Joule
Year: 1840
Title: to discover modern absolute system of electric measurements.

In 1840, he graduated his tangent galvanometer to correspond with the system of electric measurement he had adopted. The electric currents used in his experiments were thenceforth measured on the new system; and the numbers given in Joule's papers from 1840 downward are easily reducible to the modern absolute system of electric measurements.

*J. J. Nervander
Scientist’s Name: J.J. Neravander
Year: 1834
Title: to improve the measurements of electric current.

J.J. Nervander designed the more- sensitive tangent galvanometer in 1834, which led to a great improvement in precise measurements of electric current. Because of its ingenuous coiling arrangements, Germander was able to use the tangent busily to prove the validity of the law that the tangent of the deviation angle of the needle of the tangent-bus sol is proportional to the electric current flowing through its coil.

*Lord Kelvin's magneto-static tangent galvanometer
Scientist’s Name: Lord Kelvin
Year: 1887
Title: discovery of tangent galvanometer as lamp.

This form of the tangent galvanometer is designed by Lord Kelvin in c.1887. This is a magneto-static tangent galvanometer used as a lamp counter. The instrument originally consisted of a small magnet on an aluminium pointer suspended at the centre of two loops of heavy copper ribbon positioned above two sets of strong bar magnets.
It is very similar in construction to GLAHM 113325 suggesting that it was designed for use in a lighting system such as the one in Lord Kelvin's laboratory and lecture theatre.

THEORY AND WORKING

Construction
A TG consists of a circular coil of insulated copper wire wound on a circular non magnetic frame. The frame is mounted vertically on a horizontal base provided with levelling screws on the base. The coil can be rotated on a vertical axis passing through its centre. A compass box is mounted horizontally at the centre of a circular scale. The compass box is circular in shape. It consists of a tiny, powerful magnetic needle pivoted at the centre of the coil. The magnetic needle is free to rotate in the horizontal plane. The circular scale is divided into four quadrants. Each quadrant is graduated from 0° to 90°. A long thin aluminium pointer is attached to the needle at its centre and at right angle to it. To avoid errors due to parallax a plane mirror is mounted below the compass needle.

Theory
When current is passed through the TG a magnetic field is created at its centre given by where I is the current in ampere, n is the number of turns of the coil and r is the radius of the coil.
If the TG is set such that the plane of the coil is along the magnetic meridian i.e. B is perpendicular to ( is the horizontal component of the Earths magnetic field), the needle rests along the resultant. From tangent law, i.e.
or
where K is called the Reduction Factor of the TG.

Working
In the tangent galvanometer there is a circular coil having one or more turns of wire, at the centre of which a magnetic needle is either balanced on a point or suspended by a fine fibre of silk or quartz. The instrument is placed so that the plane of the coil is vertical and in the magnetic north and south plane (Figure 5(A)).

FIGURE 5(A)
When a current is sent through the coil the needle turns to one side or the other, and the strength of the current is proportional to the tangent of the angle of deflection. The force due to the current in the coil is at right angles to the plane of the coil at its centre and the strength of the field at that point in a given coil is proportional to the strength of the current (Figure 5(B)).

FIGURE 5(B)
Let G represent the strength of field at the centre due to the coil when unit current is flowing, then IG will be the strength of field when the current strength is I. Let OA in Figure 5(B) represent the plane of the coil and O the point where the needle is placed, then when no current is flowing the needle points in the direction OA, being acted on only by the horizontal component H of the earth's magnetic force. The magnetic force F due to the current in the coil is IG and at right angles to H, therefore, the resultant force R is the diagonal of the rectangle whose sides are IG and H, and

Where x is the angle which the resultant force makes with H. But the needle must point in the direction of the resultant force, and so x is the angle through which the needle turns. Therefore

And if H and G are known the current may be determined by measuring the angle x. In case of a tangent galvanometer the magnetic force F due to the coil is expressed by IG.
But if the current is measured in electromagnetic units,

And since the length of n turns of wire of radius r is,

The galvanometer coil constant G can be calculated from this formula when the coil of the galvanometer has so large a radius compared with the length of the needle that the poles of the needle may be regarded as at the centre, and when the cross section of the coil is so small that all the turns bear nearly the same relation to the needle. If G is determined in this way, r being measured in centimetres, and if H is found in C.G.S. units system, the current will be also found in C.G.S. electromagnetic units by the use of the formula:

To obtain the current strength in amperes, we must take as the value of the coil constant:

By this method the strength of a current is determined in amperes directly from the fundamental units of length, mass, and time, for we have already seen how the measurement of H is based on these units. A tangent galvanometer in which the constant is determined in this way directly from measurements of the coil is known as a standard galvanometer.

SUMMARY

Galvanometers were the first instruments used to determine the presence, direction, and strength of an electric current in a conductor. All galvanometers are based upon the discovery by Hans C. Oersted that a magnetic needle is deflected by the presence of an electric current in a nearby conductor. The extent to which the needle turns is dependent upon the strength of the current. These meters were called tangent galvanometers because the tangent of the angle of deflection of the needle is proportional to the strength of the current in the coil. A tangent galvanometer consists of a coil of insulated copper wire wound on a circular non-magnetic frame. It works on the basis of tangent law of magnetism.It works by using a compass needle to compare a magnetic field generated by the unknown current to the magnetic field of the Earth. It gets its name from its operating principle, the tangent law of magnetism, which states that the tangent of the angle a compass needle makes is proportional to the ratio of the strengths of the two perpendicular magnetic fields.

Struers Tangent Galvanometer
Unfortunately, simple galvanometers such as the Struers model shown above were inaccurate and inconsistent in their readings. By placing the compass at the centre of a precisely calculated circle, accuracy could be improved substantially (see down). Other improvements were added later including replacing the compass with a specially designed meter movement, adding levelling screws, etc.

Central Scientific Tangent Galvanometer utilizing compass (1941)

These large stationary-coil type galvanometers were used as the standard current measuring instrument into the last quarter of the 19th century. Additional examples of tangent galvanometers are shown below:

Harris Tangent Galvanometer
Harris Tangent Galvanometer
Eureka Scientific Tangent Galvanometer
University Supply Tangent Galvanometer

Knott Tangent
Galvanometer
Early Tangent Galvanometer
Early Rectangular Tangent Galvanometer
University Supply Tangent Galvanometer

Suggestion: One of the limitations of tangent galvanometers was that the length of the needle had to be kept very short in order to minimize the effects of the earth's magnetic field and reduce damping errors introduced by the mass of the needle itself. Unfortunately, the shorter the needle, the less distance the tip will travel as it inscribes an arc, and thus the more difficult it will be to read very small changes in current.

This problem is solved ingeniously by using a beam of light as the needle; a shaft is placed through the centre of the needle and a very small mirror is attached. A beam of light is reflected off of the mirror and onto a scale located about three feet away. The result is that an extremely small deflection of the mirror will cause a much larger movement of the beam on the scale. These types of galvanometers are called Reflecting Galvanometer.

BIBLIOGRAPHY

1. http://physics.kenyon.edu/EarlyApparatus/Electrical_Measurements/Tangent_Galvanometer/Tangent_Galvanometer.html
2. http://en.wikipedia.org/wiki/Galvanometer
3. http://www.historicalprintshop.com/web_pages/S/science/scientific.instruments/scientific.instruments.html
4. http://www.scran.ac.uk/database/record.php?usi=000-000-529-673-C&&
5. http://ieeexplore.ieee.org/Xplore/login.jsp?url=/iel5/5289/4534362/04534374.pdf?arnumber=4534374
6. http://www.scran.ac.uk/database/record.php?usi=000-000-529-497-C
7. http://www.economicexpert.com/a/Tangent:galvanometer.html
8. http://chem.ch.huji.ac.il/instruments/test/galvanometers.htm
9. http://www.sparkmuseum.com/GALV.HTM

optical isomerism

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:34:00 -0700

OPTICAL ISOMERISM

SUBMITTED BY- Shweta Bhardwaj
COURSE - CHE-155
PROGRAMME- BSc (hons.) BIOTECH
PROGRAMME CODE- 178
ROLL NO.- R280A03
REGISTREATION NO.- 10801595
SUBMITTED TO- Dr. Ramesh Thakur

OPTICAL ISOMERISM
Optical isomerism is a form of stereoisomerism. This page explains what stereoisomers are and how you recognise the possibility of optical isomers in a molecule.
What is stereoisomerism?
What are isomers?
Isomers are molecules that have the same molecular formula, but have a different arrangement of the atoms in space. That excludes any different arrangements which are simply due to the molecule rotating as a whole, or rotating about particular bonds.
Where the atoms making up the various isomers are joined up in a different order, this is known as structural isomerism. Structural isomerism is not a form of stereoisomerism, and is dealt with on a separate page.

What are stereoisomers?
In stereoisomerism, the atoms making up the isomers are joined up in the same order, but still manage to have a different spatial arrangement. Optical isomerism is one form of stereoisomerism.
Optical isomerism
Why optical isomers?
Optical isomers are named like this because of their effect on plane polarised light.

Simple substances which show optical isomerism exist as two isomers known as enantiomers.
• A solution of one enantiomer rotates the plane of polarisation in a clockwise direction. This enantiomer is known as the (+) form.
For example, one of the optical isomers (enantiomers) of the amino acid alanine is known as (+)alanine.
• A solution of the other enantiomer rotates the plane of polarisation in an anti-clockwise direction. This enantiomer is known as the (-) form. So the other enantiomer of alanine is known as or (-)alanine.
• If the solutions are equally concentrated the amount of rotation caused by the two isomers is exactly the same - but in opposite directions.
• When optically active substances are made in the lab, they often occur as a 50/50 mixture of the two enantiomers. This is known as a racemic mixture or racemate. It has no effect on plane polarised light.

How optical isomers arise
The examples of organic optical isomers required at A' level all contain a carbon atom joined to four different groups. These two models each have the same groups joined to the central carbon atom, but still manage to be different:

Obviously as they are drawn, the orange and blue groups aren't aligned the same way. Could you get them to align by rotating one of the molecules? The next diagram shows what happens if you rotate molecule B.

They still aren't the same - and there is no way that you can rotate them so that they look exactly the same. These are isomers of each other.
They are described as being non-superimposable in the sense that (if you imagine molecule B being turned into a ghostly version of itself) you couldn't slide one molecule exactly over the other one. Something would always be pointing in the wrong direction.

What happens if two of the groups attached to the central carbon atom are the same? The next diagram shows this possibility.

The two models are aligned exactly as before, but the orange group has been replaced by another pink one.
Rotating molecule B this time shows that it is exactly the same as molecule A. You only get optical isomers if all four groups attached to the central carbon are different.

Chiral and achiral molecules
The essential difference between the two examples we've looked at lies in the symmetry of the molecules.
If there are two groups the same attached to the central carbon atom, the molecule has a plane of symmetry. If you imagine slicing through the molecule, the left-hand side is an exact reflection of the right-hand side.
Where there are four groups attached, there is no symmetry anywhere in the molecule.

A molecule which has no plane of symmetry is described as chiral. The carbon atom with the four different groups attached which causes this lack of symmetry is described as a chiral centre or as an asymmetric carbon atom.
The molecule on the left above (with a plane of symmetry) is described as achiral.
Only chiral molecules have optical isomers.
The relationship between the enantiomers
One of the enantiomers is simply a non-superimposable mirror image of the other one.
In other words, if one isomer looked in a mirror, what it would see is the other one. The two isomers (the original one and its mirror image) have a different spatial arrangement, and so can't be superimposed on each other.

If an achiral molecule (one with a plane of symmetry) looked in a mirror, you would always find that by rotating the image in space, you could make the two look identical. It would be possible to superimpose the original molecule and its mirror image.
Some real examples of optical isomers
Butan-2-ol
The asymmetric carbon atom in a compound (the one with four different groups attached) is often shown by a star.

It's extremely important to draw the isomers correctly. Draw one of them using standard bond notation to show the 3-dimensional arrangement around the asymmetric carbon atom. Then draw the mirror to show the examiner that you know what you are doing, and then the mirror image.

Notice that you don't literally draw the mirror images of all the letters and numbers! It is, however, quite useful to reverse large groups - look, for example, at the ethyl group at the top of the diagram.
It doesn't matter in the least in what order you draw the four groups around the central carbon. As long as your mirror image is drawn accurately, you will automatically have drawn the two isomers.
So which of these two isomers is (+)butan-2-ol and which is (-)butan-2-ol? There is no simple way of telling that. For A'level purposes, you can just ignore that problem - all you need to be able to do is to draw the two isomers correctly.
2-hydroxypropanoic acid (lactic acid)
Once again the chiral centre is shown by a star.

The two enantiomers are:

It is important this time to draw the COOH group backwards in the mirror image. If you don't there is a good chance of you joining it on to the central carbon wrongly.

If you draw it like this in an exam, you won't get the mark for that isomer even if you have drawn everything else perfectly.
2-aminopropanoic acid (alanine)
This is typical of naturally-occurring amino acids. Structurally, it is just like the last example, except that the -OH group is replaced by -NH2

The two enantiomers are:

Only one of these isomers occurs naturally: the (+) form. You can't tell just by looking at the structures which this is.
It has, however, been possible to work out which of these structures is which. Naturally occurring alanine is the right-hand structure, and the way the groups are arranged around the central carbon atom is known as an L- configuration. Notice the use of the capital L. The other configuration is known as D-.
So you may well find alanine described as L-(+)alanine.
That means that it has this particular structure and rotates the plane of polarisation clockwise.
Even if you know that a different compound has an arrangement of groups similar to alanine, you still can't say which way it will rotate the plane of polarisation.
The other amino acids, for example, have the same arrangement of groups as alanine does (all that changes is the CH3 group), but some are (+) forms and others are (-) forms.
It's quite common for natural systems to only work with one of the enantiomers of an optically active substance. It isn't too difficult to see why that might be. Because the molecules have different spatial arrangements of their various groups, only one of them is likely to fit properly into the active sites on the enzymes they work with.
In the lab, it is quite common to produce equal amounts of both forms of a compound when it is synthesised. This happens just by chance, and you tend to get racemic mixtures.
________________________________________
Note: For a detailed discussion of this, you could have a look at the page on the addition of HCN to aldehydes
________________________________________
Chirality

Two enantiomers of a generic amino acid

The two optical isomers of alanine.

The two enantiomers of bromochlorofluoromethane
The term chiral (pronounced /ˈkaɪrəl/) is used to describe an object that is non-superposable on its mirror image.
Human hands are perhaps the most universally recognized example of chirality: The left hand is a non-superposable mirror image of the right hand; no matter how the two hands are oriented, it is impossible for all the major features of both hands to coincide. This difference in symmetry becomes obvious if someone attempts to shake the right hand of a person using his left hand, or if a left-handed glove is placed on a right hand. The term chirality is derived from the Greek word for hand, χειρ (/cheir/).
When used in the context of chemistry, chirality usually refers to molecules. Two mirror images of a molecule that cannot be superposed onto each other are referred to as enantiomers or optical isomers. Because the difference between right and left hands is universally known and easy to observe, many pairs of enantiomers are designated as "right-" and "left-handed." A mixture of equal amounts of the two enantiomers is said to be a racemic mixture. Molecular chirality is of interest because of its application to stereochemistry in inorganic chemistry, organic chemistry, physical chemistry, biochemistry, and supramolecular chemistry.
The symmetry of a molecule (or any other object) determines whether it is chiral. A molecule is achiral (not chiral) if and only if it has an axis of improper rotation; that is, an n-fold rotation (rotation by 360°/n) followed by a reflection in the plane perpendicular to this axis that maps the molecule onto itself. (See chirality (mathematics).) A simplified rule applies to tetrahedrally-bonded carbon, as shown in the illustration: if all four substituents are different, the molecule is chiral. A chiral molecule is not necessarily asymmetric, that is, devoid of any symmetry elements, as it can have, for example, rotational symmetry.

History
The term optical activity is derived from the interaction of chiral materials with polarized light. A solution of the (−)-form of an optical isomer rotates the plane of polarization of a beam of plane polarized light in a counterclockwise direction, vice-versa for the (+) optical isomer. The property was first observed by Jean-Baptiste Biot in 1815 [1], and gained considerable importance in the sugar industry, analytical chemistry, and pharmaceuticals. Louis Pasteur deduced in 1848 that this phenomenon has a molecular basis[2]. Artificial composite materials displaying the analog of optical activity but in the microwave region were introduced by J.C. Bose in 1898 [3], and gained considerable attention from the mid-1980s [4]. The term chirality itself was coined by Lord Kelvin in 1873.[1]
The word “racemic” is derived from the Latin word for grape; the term having its origins in the work of Louis Pasteur who isolated racemic tartaric acid from wine.
Naming conventions
By configuration: R- and S-
For chemists, the R / S system is the most important nomenclature system for denoting enantiomers, which does not involve a reference molecule such as glyceraldehyde. It labels each chiral center R or S according to a system by which its substituents are each assigned a priority, according to the Cahn Ingold Prelog priority rules(CIP), based on atomic number. If the center is oriented so that the lowest-priority of the four is pointed away from a viewer, the viewer will then see two possibilities: If the priority of the remaining three substituents decreases in clockwise direction, it is labeled R (for Rectus), if it decreases in counterclockwise direction, it is S (for Sinister).
This system labels each chiral center in a molecule (and also has an extension to chiral molecules not involving chiral centers). Thus, it has greater generality than the D/L system, and can label, for example, an (R,R) isomer versus an (R,S) — diastereomers.
The R / S system has no fixed relation to the (+)/(−) system. An R isomer can be either dextrorotatory or levorotatory, depending on its exact substituents.
The R / S system also has no fixed relation to the D/L system. For example, the side-chain one of serine contains a hydroxyl group, -OH. If a thiol group, -SH, were swapped in for it, the D/L labeling would, by its definition, not be affected by the substitution. But this substitution would invert the molecule's R / S labeling, because the CIP priority of CH2OH is lower than that for CO2H but the CIP priority of CH2SH is higher than that for CO2H.
For this reason, the D/L system remains in common use in certain areas of biochemistry, such as amino acid and carbohydrate chemistry, because it is convenient to have the same chiral label for all of the commonly occurring structures of a given type of structure in higher organisms. In the D/L system, they are all L; in the R / S system, they are mostly S but there are some common exceptions.
By optical activity: (+)- and (−)-
An enantiomer can be named by the direction in which it rotates the plane of polarized light. If it rotates the light clockwise (as seen by a viewer towards whom the light is traveling), that enantiomer is labeled (+). Its mirror-image is labeled (−). The (+) and (−) isomers have also been termed d- and l-, respectively (for dextrorotatory and levorotatory). This labeling is easy to confuse with D- and L-.
By configuration: D- and L-
An optical isomer can be named by the spatial configuration of its atoms. The D/L system does this by relating the molecule to glyceraldehyde. Glyceraldehyde is chiral itself, and its two isomers are labeled D and L. Certain chemical manipulations can be performed on glyceraldehyde without affecting its configuration, and its historical use for this purpose (possibly combined with its convenience as one of the smallest commonly used chiral molecules) has resulted in its use for nomenclature. In this system, compounds are named by analogy to glyceraldehyde, which, in general, produces unambiguous designations, but is easiest to see in the small biomolecules similar to glyceraldehyde. One example is the amino acid alanine, which has two optical isomers, and they are labeled according to which isomer of glyceraldehyde they come from. On the other hand, glycine, the amino acid derived from glyceraldehyde, has no optical activity, as it is not chiral (achiral). Alanine, however, is chiral.
The D/L labeling is unrelated to (+)/(−); it does not indicate which enantiomer is dextrorotatory and which is levorotatory. Rather, it says that the compound's stereochemistry is related to that of the dextrorotatory or levorotatory enantiomer of glyceraldehyde—the dextrorotatory isomer of glyceraldehyde is, in fact, the D isomer. Nine of the nineteen L-amino acids commonly found in proteins are dextrorotatory (at a wavelength of 589 nm), and D-fructose is also referred to as levulose because it is levorotatory.
A rule of thumb for determining the D/L isomeric form of an amino acid is the "CORN" rule. The groups:
COOH, R, NH2 and H (where R is a variant carbon chain)
are arranged around the chiral center carbon atom. Sighting with the hydrogen atom away from the viewer, if these groups are arranged clockwise around the carbon atom, then it is the D-form. If counter-clockwise, it is the L-form.
Nomenclature
• Any non-racemic chiral substance is called scalemic [2]
• A chiral substance is enantiopure or homochiral when only one of two possible enantiomers is present.
• A chiral substance is enantioenriched or heterochiral when an excess of one enantiomer is present but not to the exclusion of the other.
• Enantiomeric excess or ee is a measure for how much of one enantiomer is present compared to the other. For example, in a sample with 40% ee in R, the remaining 60% is racemic with 30% of R and 30% of S, so that the total amount of R is 70%.
Types
In general, chiral molecules have point chirality, centering around a single atom, usually carbon, which has four different substituents. The two enantiomers of such compounds are said to have different absolute configurations at this center. This center is thus stereogenic (i.e., a grouping within a molecular entity that may be considered a focus of stereoisomerism), and is exemplified by the α-carbon of amino acids. A molecule can have multiple chiral centers without being chiral overall if there is a symmetry element (a mirror plane or inversion center), which relates the two (or more) chiral centers. Such a molecule is called a meso compound. It is also possible for a molecule to be chiral without having actual point chirality. Common examples include 1,1'-bi-2-naphthol (BINOL) and 1,3-dichloro-allene, which have axial chirality, and (E)-cyclooctene, which has planar chirality.
It is important to keep in mind that molecules that are dissolved in solution or are in the gas phase usually have considerable flexibility, and, thus, may adopt a variety of different conformations. These various conformations are themselves almost always chiral. However, when assessing chirality, one must use a structural picture of the molecule that corresponds to just one chemical conformation - the most symmetric conformation possible.
When the optical rotation for an enantiomer is too low for practical measurement it is said to exhibit cryptochirality.
Even isotopic differences must be considered when examining chirality. Replacing one of the two 1H atoms at the CH2 position of benzyl alcohol with a deuterium (²H) makes that carbon a stereocenter. The resulting benzyl-α-d alcohol exists as two distinct enantiomers, which can be assigned by the usual stereochemical naming conventions. The S enantiomer has [α]D = +0.715°.[5]
Properties of enantiomers
Enantiomers are identical with respect to ordinary chemical reactions and properties (i.e., will have identical Rfs by TLC, identical NMR spectra, identical IR spectra), but differences arise when they are in the presence of other chiral molecules or objects. Different enantiomers of chiral compounds often taste and smell differently and have different effects as drugs - see below.
One chiral 'object' that interacts differently with the two enantiomers of a chiral compound is circularly polarised light: An enantiomer will absorb left- and right-circularly polarised light to differing degrees. This is the basis of circular dichroism (CD) spectroscopy. Usually the difference in absorptivity is relatively small (parts per thousand). CD spectroscopy is a powerful analytical technique for investigating the secondary structure of proteins and for determining the absolute configurations of chiral compounds, in particular, transition metal complexes. CD spectroscopy is replacing polarimetry as a method for characterising chiral compounds, although the latter is still popular with sugar chemists.
In biology
Many biologically active molecules are chiral, including the naturally occurring amino acids (the building blocks of proteins), and sugars. In biological systems, most of these compounds are of the same chirality: most amino acids are L and sugars are D. Typical naturally occurring proteins, made of L amino acids, are known as left-handed proteins, whereas D amino acids produce right-handed proteins.
The origin of this homochirality in biology is the subject of much debate.[6] Most scientists believe that Earth life's “choice” of chirality was purely random, and that if carbon-based life forms exist elsewhere in the universe, their chemistry could theoretically have opposite chirality.
Enzymes, which are chiral, often distinguish between the two enantiomers of a chiral substrate. Imagine an enzyme as having a glove-like cavity that binds a substrate. If this glove is right-handed, then one enantiomer will fit inside and be bound, whereas the other enantiomer will have a poor fit and is unlikely to bind.
D-form amino acids tend to taste sweet, whereas L-forms are usually tasteless. Spearmint leaves and caraway seeds, respectively, contain L-carvone and D-carvone - enantiomers of carvone. These smell different to most people because our olfactory receptors also contain chiral molecules that behave differently in the presence of different enantiomers.
Chirality is important in context of ordered phases as well, for example the addition of a small amount of an optically active molecule to a nematic phase (a phase that has long range orientational order of molecules) transforms that phase to a chiral nematic phase (or cholesteric phase). Chirality in context of such phases in polymeric fluids has also been studied in this context (Srinivasarao, 1999).
In drugs
Many chiral drugs must be made with high enantiomeric purity due to potential side-effects of the other enantiomer. (The other enantiomer may also merely be inactive.)
• Thalidomide: Thalidomide is racemic. One enantiomer is effective against morning sickness, whereas the other is teratogenic. In this case, administering just one of the enantiomers to a pregnant patient does not help, as the two enantiomers are readily interconverted in vivo. Thus, if a person is given either enantiomer, both the D and L isomers will eventually be present in the patient's serum.
• Ethambutol: Whereas one enantiomer is used to treat tuberculosis, the other causes blindness.
• Naproxen: One enantiomer is used to treat arthritis pain, but the other causes liver poisoning with no analgesic effect.
• Steroid receptor sites also show stereoisomer specificity.
• Penicillin's activity is stereodependent. The antibiotic must mimic the D-alanine chains that occur in the cell walls of bacteria in order to react with and subsequently inhibit bacterial transpeptidase enzyme.
• Only L-propranolol is a powerful adrenoceptor antagonist, whereas D-propranolol is not. However, both have local anesthetic effect.
• The L-isomer of Methorphan, levomethorphan is a potent opioid analgesic, while the D-isomer, dextromethorphan is a dissociative cough suppressant.
• S(-) isomer of carvedilol, a drug that interacts with adrenoceptors, is 100 times more potent as beta receptor blocker than R(+) isomer. However, both the isomers are approximately equipotent as alpha receptor blockers.
• The D-isomers of amphetamine and methamphetamine are strong CNS stimulants, while the L-isomers of both drugs lack appreciable CNS(central nervous system) stimulant effects, but instead stimulate the peripheral nervous system. For this reason, the Levo-isomer of methamphetamine is available as an OTC nasal inhaler in some countries, while the Dextro-isomer is banned from medical use in all but a few countries in the world, and highly regulated in those countries who do allow it to be used medically.
In inorganic chemistry
Many coordination compounds are chiral; for example, the well-known [Ru(2,2'-bipyridine)3]2+ complex in which the three bipyridine ligands adopt a chiral propeller-like arrangement [7]. In this case, the Ru atom may be regarded as a stereogenic center, with the complex having point chirality. The two enantiomers of complexes such as [Ru(2,2'-bipyridine)3]2+ may be designated as Λ (left-handed twist of the propeller described by the ligands) and Δ (right-handed twist). Hexol is a chiral cobalt complex that was first investigated by Alfred Werner. Resolved hexol is significant as being the first compound devoid of carbon to display optical activity.
Chirality of amines

Tertiary amines (see image) are chiral in a way similar to carbon compounds: The nitrogen atom bears four distinct substituents counting the lone pair. However, the energy barrier for the inversion of the stereocenter is, in general, about 30 kJ/mol, which means that the two stereoisomers are rapidly interconverted at room temperature. As a result, amines such as NHRR' cannot be resolved optically and NRR'R" can only be resolved when the R, R', and R" groups are constrained in cyclic structures.
Theory of origin
A paper published in February 29, 2008 by researchers led by Sandra Pizzarello, from Arizona State University, reveals that the Murchison meteorite contains sizable molecular asymmetry of up to 14%, "giving support to the idea that biomolecular traits such as chiral asymmetry could have been seeded in abiotic chemistry ahead of life."[8]
"Thanks to the pristine nature of this meteorite, we were able to demonstrate that other extraterrestrial amino acids carry the left-handed excesses in meteorites and, above all, that these excesses appear to signify that their precursor molecules, the aldehydes, also carried such excesses," Pizzarello said. "In other words, a molecular trait that defines life seems to have broader distribution as well as a long cosmic lineage."[3]
Other theories of the origin of chirality on Earth have also been proposed, such as the weak nuclear force.[4]
Chemical chirality in Fiction
Although little was known about chemical chirality in the time of Lewis Carroll, his work Through the Looking-glass contains a prescient reference to the differing biological activities of enantiomeric drugs: "Perhaps Looking-glass milk isn't good to drink," Alice said to her cat.
In James Blish's Star Trek novella Spock Must Die! the tachyon 'mirrored' Mr Spock is later discovered to have stolen chemical reagents from the medical bay and to have been using them to convert certain amino acids to opposite-chirality isomers, since the mirrored Mr Spock's metabolism is reversed, and, hence, must process the opposite polarity of these isomers.
In Larry Niven's Destiny's Road, the title planet's indigenous life is based upon right-handed proteins. When human colonists arrive from Earth via a generation ship, extreme measures are taken to permit the colony's survival. A peninsula is sterilized with a lander's fusion drive, creating the titular "road" out of fused bedrock. The area is then reseeded with Earth life to provide the colonists with food. Though the soil lacks potassium due to other factors, necessitating supplements that produce a hydraulic empire common to Niven's fiction, the colony otherwise prospers. Native viruses and bacteria cannot infect colonists, resulting in longer lifespans. Sealife quickly recovers, and is consumed by the colonists as a "diet" food, as their digestive systems cannot metabolize it into fat.
In the Trauma Center series of games, doctors test for a "chiral reaction" in order to determine whether or not a patient is infected with "Gangliated Utrophin Immuno Latency Toxin," a fictional, parasitic pathogen more commonly referred to as G.U.I.L.T. A positive reaction means the patient is infected, while a negative reaction means the patient has either been cured or is not infected.
translation of French original, published by Alembic Club Reprints (Vol. 14, pp. 1-46) in 1905, facsimile reproduction by SPIE in a 1990 book". Notes
Structure and Optical Isomerism
A very important feature of the structure of amino acids (and other kinds of compounds as well, for that matter) is called optical isomerism. It applies to all amino acids except glycine.
Look at the number-two carbon atom. You should notice that in one direction it is bonded to an amino group. In another direction, it is bonded to a carboxylic group. It is also bonded to a hydrogen atom and an alkyl group or some other kind of group. Except in the case of glycine where -R is a -H, that number two carbon atom is bonded to four different groups. A carbon atom which is attached to four different groups is called an asymmetric carbon atom or sometimes a chiral carbon atom. The importance of this depends on some structural properties that we will investigate in this section.

If you are in the lab you get a model kit and follow along with the diagrams shown here. Get a carbon atom and attach to it four different groups. For convenience just use different colored units, rather than actually building an amino group and a carboxylic acid group and an isopropyl group or something like that. Then make the other models as they are shown bleow. If you are not in the lab now, you should work with the models to do this exercise when you are in the lab.
Here is a model of a carbon atom with four different groups attached.

Here is another model constructed to be the mirror image of the first model. To do this, construct a model that would appear just as the first model that you made would look like if you were looking at it in a mirror.

Here you can see why these are called mirror images of one another.

We can demonstrate that these two structures are not identical to one another by trying to superimpose one structure on another and get all of the same colored units to be in the identical places. You can see that is not possible.

The two structures are different. They are isomers of one another. It so happens that they are called optical isomers of one another because they have optical properties that are different from one another. We will discuss that particular property a little bit more when we discuss carbohydrates in a later lesson.

When asymmetric carbon atoms are present in a molecular compound, there are two ways in which the groups attached to that carbon can be arranged in the three dimensions, as we have just shown with the two models above. It is generally true, if not universally true, that only one of these optical isomers is biologically active. In other words, when these compounds are made by a plant or animal, only one of the two forms is made. When it comes time for these molecules to interact with an enzyme, only one of these molecules would react. The other would not. Both shape and orientation in biological compounds are extremely important.
Chemically, optical isomers behave the same. Biologically, they do not. One will react properly, but the other will not. Optically, there is also that difference which will be pointed out when we deal with carbohydrates in a later lesson.
We can use these models to illustrate why you need to have four different groups bonded to the central atom. One group (the black group) has been removed from the model on the left and replaced it with a duplicate of one of the other three groups (the white group). We now have a model with the central atom bonded to four groups, but they are not all different. The same has been done to the mirror image (unfortunately, you cannot see that).

By turning the second model in the right way you can see that it is identical to the first one.

Consequently, this central atom is not an asymmetric carbon atom, the molecule is not an optically active molecule, and these are identicalcompounds and not optical isomers.

References
1. ^ Pedro Cintas. "Tracing the Origins and Evolution of Chirality and Handedness in Chemical Language". Angewandte Chemie International Edition 46 (22): 4016-4024. doi:10.1002/anie.200603714.
2. ^ Infelicitous stereochemical nomenclatures for stereochemical nomenclature
3. ^ Arizona State University (2008, February 29). Key To Life Before Its Origin On Earth May Have Been Discovered. ScienceDaily. Retrieved June 16, 2008, from http://www.sciencedaily.com/releases/2008/02/080228174823.htm
4. ^ Castelvecchi, Davide. (2007). Alien Pizza, Anyone?, Science News vol. 172, pp. 107-109. (references

Hydrogen fuel cell as way out for energy crises, basic technology used and latest advances, applications

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:31:00 -0700

Hydrogen fuel cell as way out for energy crises, basic technology used and latest advances, applications
Acknowledgement

Gratitude cannot be seen or expressed. It can only be felt in heart and is beyond description. Often words are inadequate to serve as a model of expression of one’s feeling, specially the sense of indebtedness and gratitude to all those who help us in our duty.

It is of immense pleasure and profound privilege to express my gratitude and indebtedness along with sincere thanks to Dr Kailash Juglan, lecturer of PHYSICS of Lovely Professional University for providing me the opportunity to work for a project on “Hydrogen fuel cell as way out for energy crises, basic technology used and latest advances, applications ”

I am beholden to my family and friends for their blessings and encouragement.

Always Obediently
Prateek Joshi

What Is Energy Crisis?

An energy crisis is any great bottleneck (or price rise) in the supply of energy resources to an economy. It usually refers to the shortage of oil and additionally to electricity or other natural resources. An energy crisis may be referred to as an oil crisis, petroleum crisis, energy shortage, electricity shortage or electricity crisis

Market failure is possible when monopoly manipulation of markets occurs. A crisis can develop due to industrial actions like union organized strikes and government embargoes. The cause may be over-consumption, ageing infrastructure, choke point disruption or bottlenecks at oil refineries and port facilities that restrict fuel supply. An emergency may emerge during unusually cold winters.this probabaly rises the depletion of energy.
Pipeline failures and other accidents may cause minor interruptions to energy supplies. A crisis could possibly emerge after infrastructure damage from severe weather. Attacks by terrorists or militia on important infrastructure are a possible problem for energy consumers, with a successful strike on a Middle East facility potentially causing global shortages. Political events, for example, when governments change due to regime change, monarchy collapse, military occupation, and coup may disrupt oil and gas production and create shortages.

Energy Crisis in History

•
• 1973 oil crisis - Cause: an OPEC oil export embargo by many of the major Arab oil-producing states, in response to western support of Israel during the Yom Kippur War
• 1979 energy crisis - Cause: the Iranian revolution
• 1990 spike in the price of oil Cause: the Gulf War
• The 2000–2001 California electricity crisis - Cause: failed deregulation, and business corruption.
• The UK fuel protest of 2000 - Cause: Raise in the price of crude oil combined with already relatively high taxation on road fuel in the UK.
• North American Gas crisis
• Argentine gas crisis of 2004
• North Korea has had energy shortages for many years.
• Zimbabwe has experienced a shortage of energy supplies for many years due to financial mismanagement.
While not entering a full crisis, political riots that occurred during the 2007 Burmese anti-government protests were initially sparked by rising energy prices. Likewise the Russia-Ukraine gas dispute and the Russia-Belarus energy dispute have been mostly resolved before entering a prolonged crisis stage.

Present Day Crisis

Crises that currently exist include:

• Oil Price Increases in 2003 - Caused by continued global increases in petroleum demand coupled with production stagnation, the falling value of the U.S. dollar
• 2008 Central Asia energy crisis, caused by abnormally cold temperatures and low water levels in an area dependent on hydroelectric power. Despite having significant hydrocarbon reserves, in February 2008 the President of Pakistan announced plans to tackle energy shortages that were reaching crisis stage. At the same time the South African President was appeasing fears of a prolonged electricity crisis in South Africa.
• South African electrical crisis. The South African crisis, which may last to 2012, lead to large price rises for platinum in February 2008 and reduced gold production.
• China experienced severe energy shortages towards the end of 2005 and again in early 2008. During the latter crisis they suffered severe damage to power networks along with diesel and coal shortages. Supplies of electricity in Guangdong province, the manufacturing hub of China, are predicted to fall short by an estimated 10 GW

Predictions

Although technology has made oil extraction more efficient, the world is having to struggle to provide oil by using increasingly costly and less productive methods such as deep sea drilling, and developing environmentally sensitive areas such as the Arctic National Wildlife Refuge.
The world's population continues to grow at a quarter of a million people per day, increasing the consumption of energy. Although far less from people in developing countries, especially USA, the per capita energy consumption of China, India and other developing nations continues to increase as the people living in these countries adopt more energy intensive lifestyles. At present a small part of the world's population consumes a large part of its resources, with the United States and its population of 300 million people consuming far more oil than China with its population of 1.3 billion people.

Future and alternative energy sources

In response to the petroleum crisis, the principles of green energy and sustainable living movements gain popularity. This has led to increasing interest in alternate power/fuel research such as fuel cell technology, liquid nitrogen economy, hydrogen fuel, biomethanol, biodiesel, Karrick process, solar energy, geothermal energy, tidal energy, wave power, and wind energy, and fusio power. To date, only hydroelectricity and nuclear power have been significant alternatives to fossil fuel.
Hydrogen gas is currently produced at a net energy loss from natural gas, which is also experiencing declining production in North America and elsewhere. When not produced from natural gas, hydrogen still needs another source of energy to create it, also at a loss during the process. This has led to hydrogen being regarded as a 'carrier' of energy, like electricity, rather than a 'source'. The unproven dehydrogenating process has also been suggested for the use water as an energy source.
Efficiency mechanisms such as Negawatt power can encourage significantly more effective use of current generating capacity. It is a term used to describe the trading of increased efficiency, using consumption efficiency to increase available market supply rather than by increasing plant generation capacity. As such, it is a demand-side as opposed to a supply-side measure.

Growing demand for a new fuel

AS the increase with the energy crisis over the present years there has been a growing demand for an alternative sources of energy. Some of them are

1-SolarEnergy
2-Tidal Energy
3-Hydro energy
4-Biological Energy
5-Hydrogen Energy

Much has been said about all the other form of energies except that of HYDROGEN based energy. Which has created awareness among the mass and scientific world in the past few years. As compared to any other form of energy Hydrogen based energy has two main benefits
Firstly it does not leave any residue after combustion except pure water which can be used for infinite purposes and second is that It produces a large amount of energy when its combustion takes place.

What Is Hydrogen?

Hydrogen is the chemical element with atomic number 1. It is represented by the symbol H. At standard temperature and pressure hydrogen is a colorless, odorless, nonmetallic, tasteless, highly flammable diatomic gas with the molecular formula H2. With an atomic weight of 1.00794, hydrogen is the lightest element.
Hydrogen is the most abundant of the chemical elements, constituting roughly 75% of the universe's elemental mass. Stars in the main sequence are mainly composed of hydrogen in its plasma state. Elemental hydrogen is relatively rare on Earth, and is industrially produced from hydrocarbons such as methane, after which most elemental hydrogen is used "captively" (meaning locally at the production site), with the largest markets about equally divided between fossil fuel upgrading (e.g., hydrocracking) and ammonia production (mostly for the fertilizer market). Hydrogen may be produced from water using the process of electrolysis, but this process is presently significantly more expensive commercially than hydrogen production from natural gas.
The most common naturally occurring isotope of hydrogen, known as protium, has a single proton and no neutrons. In ionic compounds it can take on either a positive charge (becoming a cation composed of a bare proton) or a negative charge (becoming an anion known as a hydride). Hydrogen can form compounds with most elements and is present in water and most organic compounds. It plays a particularly important role in acid-base chemistry, in which many reactions involve the exchange of protons between soluble molecules. As the only neutral atom for which the Schrödinger equation can be solved analytically, study of the energetics and bonding of the hydrogen atom has played a key role in the development of quantum mechanics.
The solubility and characteristics of hydrogen with various metals are very important in metallurgy (as many metals can suffer hydrogen embrittlement) and in developing safe ways to store it for use as a fuel. Hydrogen is highly soluble in many compounds composed of rare earth metals and transition metals and can be dissolved in both crystalline and amorphous metals. Hydrogen solubility in metals is influenced by local distortions or impurities in the metal crystal lattice.

History

Hydrogen gas, H2, was first artificially produced and formally described by T. Von Hohenheim (also known as Paracelsus, 1493–1541) via the mixing of metals with strong acids He was unaware that the flammable gas produced by this chemical reaction was a new chemical element. In 1671, Robert Boyle rediscovered and described the reaction between iron filings and dilute acids, which results in the production of hydrogen gas. In 1766, Henry Cavendish was the first to recognize hydrogen gas as a discrete substance, by identifying the gas from a metal-acid reaction as "inflammable air" and further finding in 1781 that the gas produces water when burned. He is usually given credit for its discovery as an element. In 1783, Antoine Lavoisier gave the element the name of hydrogen (from the Greek hydro meaning water and genes meaning creator) when he and Laplace reproduced Cavendish's finding that water is produced when hydrogen is burned
Hydrogen was liquefied for the first time by James Dewar in 1898 by using regenerative cooling and his invention, the vacuum flask. He produced solid hydrogen the next year. Deuterium was discovered in December 1931 by Harold Urey, and tritium was prepared in 1934 by Ernest Rutherford Mark Oliphant, and Paul Harteck. Heavy water, which consists of deuterium in the place of regular hydrogen, was discovered by Urey's group in 1932. François Isaac de Rivaz built the first internal combustion engine powered by a mixture of hydrogen and oxygen in 1806. Edward Daniel Clarke invented the hydrogen gas blowpipe in 1819. The Döbereiner's lamp and limelight were invented in 1823
The first hydrogen-filled balloon was invented by Jacques Charlesin 1783. Hydrogen provided the lift for the first reliable form of air-travel following the 1852 invention of the first hydrogen-lifted airship by Henri Giffard German count Ferdinand von Zeppelin promoted the idea of rigid airships lifted by hydrogen that later were called Zeppelins; the first of which had its maiden flight in 1900. Regularly-scheduled flights started in 1910 and by the outbreak of World War I in August 1914 they had carried 35,000 passengers without a serious incident. Hydrogen-lifted airships were used as observation platforms and bombers during the war.
The first non-stop transatlantic crossing was made by the British airship R34 in 1919. Regular passenger service resumed in the 1920s and the discovery of helium reserves in the United States promised increased safety, but the U.S. government refused to sell the gas for this purpose. Therefore, H2 was used in the Hindenburg airship, which was destroyed in a midair fire over New Jersey on May 6, 1937. The incident was broadcast live on radio and filmed. Ignition of leaking hydrogen as widely assumed to be the cause but later investigations pointed to ignition of the aluminized fabric coating by static electricity. But the damage to hydrogen's reputation as a lifting gas was already done.

Hydrogen As a FUEL

"President Bush’s remarks in his State-of-the-Union message proposing a big jump in funding for hydrogen and fuel cell research and development are terrific news. It’s imperative that Congress follows through now and makes available those funds. Aside from the tangible benefits of spending more on an environmentally benign area of energy that for too long has been treated - often condescendingly - like a poor orphan, the political message is of supreme significance. For decades, supporters of hydrogen and other alternative energy fields have argued until they were blue in the face, that the key ingredient missing in moving forward is national political will. President Bush’s support provides a large measure of that political will." --Peter Hoffmann, 31 January 2003 about the book: Hydrogen is the quintessential eco-fuel. This invisible, tasteless gas is the most abundant element in the universe. It is the basic building block and fuel of stars and an essential raw material in innumerable biological and chemical processes. As a completely nonpolluting fuel, it may hold the answer to growing environmental concerns about atmospheric accumulation of carbon dioxide and the resultant Greenhouse Effect. In this book Peter Hoffmann describes current research toward a hydrogen-based economy. He presents the history of hydrogen energy and discusses the environmental dangers of continued dependence on fossil fuels. Hydrogen is not an energy source but a carrier that, like electricity, must be manufactured. Today hydrogen is manufactured by "decarburizing" fossil fuels. In the future it will be derived from water and solar energy and perhaps from "cleaner" versions of nuclear energy. Because it can be made by a variety of methods, Hoffmann argues, it can be easily adapted by different countries and economies. Hoffmann acknowledges the social, political, and economic difficulties in replacing current energy systems with an entirely new one. Although the process of converting to a hydrogen-based economy would be complex, he demonstrates that the environmental and health benefits would far outweigh the costs.

Small whispers of hydrogen energy's vast potential have been heard along the fringes of industry since the oil shocks of the 1970s, but only last year did a steady drumbeat begin in the capital markets of Wall Street, Europe, and Asia. First BMW and Daimler-Chrysler, and then Ford, Honda, Toyota, GM, and others laid claim to hydrogen fuel and to the fuel cell as a new prime mover for the automobile.
An informed public may be all that is required to bring an end to the climate-destabilizing fossil era. Until this summer, though, we had no recent book on the emerging world hydrogen economy Information was available only to readers of periodicals like Peter Hoffmann's Hydrogen and Fuel Cell Letter and The International Journal of Hydrogen Energy.
Finally in August, two books. Hoffmann's chronicles hydrogen science and technology from the earliest days. Embedded in its historical narrative are explanations of these technologies and their advantages and drawbacks. He addresses the questions people are starting to ask: Why a hydrogen economy? How do you get hydrogen? What will it cost? Is it safe? Will it reduce global warming? What is its connection with solar and wind energy? The book's main drawback is the index, which is missing essential entries such as pipelines, carbon dioxide, leakage, sequestration, biomass, and embrittlement. But at last we now have a book we can use to understand the elements of this epic change.
Seth Dunn's Worldwatch Paper speaks from the environmental perspective and describes present practices with an eye to the future. He reports on a range of studies by government agencies, NGOs, universities, and corporations, all attempting to illuminate potential paths for the emerging hydrogen economy He compares this moment in the hydrogen fuel revolution to the early automobile era, which saw fierce competition among technologies before the gasoline-powered internal combustion engine won out as the standard.--Ty Cashman
"Decarbonization is just what it sounds like: taking the carbon out of hydrocarbon fuels. What is left is, of course, hydrogen. Decarbonization will be the industrial end-game strategy of a trend first detected by Cesare Marchetti in the 1970s, when he described a gradual shift, over centuries, from hydrocarbon fuels with high carbon and low hydrogen content (wood, peat, coal) to fuels with increasingly less carbon and more hydrogen (oil, natural gas), culminating, seemingly inevitably, in pure hydrogen as the principal energy carrier of an advanced industrial society.
"If hydrogen is ever to replace natural gas as a utility fuel, very large quantities obviously will have to be stored somewhere. Storage, to maintain a buffer for seasonal, daily, and hourly swings in demand, is essential with any system for the transmission of a gas. Storage facilities even out the ups and downs of demand, including temporary interruptions and breakdowns, and still permit steady, maximum-efficiency production.
"It has been suggested that huge amounts of hydrogen could be stored underground in exhausted natural gas fields, in natural or manmade caverns, or in aquifers.... The natural gas industry has long been using depleted gas and oil fields to store huge amounts of natural gas. Aquifers are similar to natural gas and oil fields in that they are porous geological formations, but without the fossil-fuel or natural gas content. Many of them feature a "caprock" formation, a layer on top of the formation that is usually saturated with water. This layer acts as a seal to keep the gas from leaking out; it works for both natural gas and the lighter hydrogen.

What Is Fuel Cell?

A fuel cell is an electrochemical conversion device. It produces electricity from fuel (on the anode side) and an oxidant (on the cathode side), which react in the presence of an electrolyte. The reactants flow into the cell, and the reaction products flow out of it, while the electrolyte remains within it. Fuel cells can operate virtually continuously as long as the necessary flows are maintained.
Fuel cells are different from electrochemical cell batteries in that they consume reactant, which must be replenished, whereas batteries store electrical energy chemically in a closed system. Additionally, while the electrodes within a battery react and change as a battery is charged or discharged, a fuel cell's electrodes are catalytic and relatively stable.
Many combinations of fuel and oxidant are possible. A hydrogen cell uses hydrogen as fuel and oxygen (usually from air) as oxidant. Other fuels include hydrocarbons and alcohols. Other oxidants include air, chlorine and chlorine dioxide.

Design And Working Of a Fuel Cell

A fuel cell works by catalysis, separating the component electrons and protons of the reactant fuel, and forcing the electrons to travel though a circuit, hence converting them to electrical power. The catalyst typically comprises a platinum group metal or alloy. Another catalytic process takes the electrons back in, combining them with the protons and the oxidant to form waste products (typically simple compounds like water and carbon dioxide).
In the archetypal hydrogen–oxygen proton exchange membrane fuel cell (PEMFC) design, a proton-conducting polymer membrane, (the electrolyte), separates the anode and cathode sides. This was called a "solid polymer electrolyte fuel cell" (SPEFC) in the early 1970s, before the proton exchange mechanism was well-understood. (Notice that "polymer electrolyte membrane" and "proton exchange membrane" result in the same acronym.)
On the anode side, hydrogen diffuses to the anode catalyst where it later dissociates into protons and electrons. These protons often react with oxidants causing them to become what is commonly referred to as multi-facilitated proton membranes (MFPM). The protons are conducted through the membrane to the cathode, but the electrons are forced to travel in an external circuit (supplying power) because the membrane is electrically insulating. On the cathode catalyst, oxygen molecules react with the electrons (which have traveled through the external circuit) and protons to form water — in this example, the only waste product, either liquid or vapor.
In addition to this pure hydrogen type, there are hydrocarbon fuels for fuel cells, including diesel, methanol (see: direct-methanol fuel cells and indirect methanol fuel cells) and chemical hydrides. The waste products with these types of fuel are carbon dioxide and water.

Construction of a low temperature PEMFC: Bipolar plate as electrode with in-milled gas channel structure, fabricated from conductive plastics (enhanced with carbon nanotubes for more conductivity); Porous carbon papers; reactive layer, usually on the polymer membrane applied; polymer membrane.

Condensation of water produced by a PEMFC on the air channel wall. The gold wire around the cell ensures the collection of electric current.
The materials used in fuel cells differ by type. In a typical membrane electrode assembly (MEA), the electrode–bipolar plates are usually made of metal, nickel or carbon nanotubes, and are coated with a catalyst (like platinum, nano iron powders or palladium) for higher efficiency. Carbon paper separates them from the electrolyte. The electrolyte could be ceramic or a membrane.
A typical PEM fuel cell produces a voltage from 0.6 V to 0.7 V at full rated load. Voltage decreases as current increases, due to several factors:
• Activation loss
• Ohmic loss (voltage drop due to resistance of the cell components and interconnects)
• Mass transport loss (depletion of reactants at catalyst sites under high loads, causing rapid loss of voltage)[3]
To deliver the desired amount of energy, the fuel cells can be combined in series and parallel circuits, where series yield higher voltage, and parallel allows a stronger current to be drawn. Such a design is called a fuel cell stack. Further, the cell surface area can be increased, to allow stronger current from each cell.

History

The principle of the fuel cell was discovered by German scientist Christian Friedrich Schönbein in 1838 and published in one of the scientific magazines of the time. Based on this work, the first fuel cell was demonstrated by Welsh scientist Sir William Robert Grove in the February 1839 edition of the Philosophical Magazine and Journal of Science and later sketched, in 1842, in the same journal. The fuel cell he made used similar materials to today's phosphoric-acid fuel cell.
In 1955, W. Thomas Grubb, a chemist working for the General Electric Company (GE), further modified the original fuel cell design by using a sulphonated polystyrene ion-exchange membrane as the electrolyte. Three years later another GE chemist, Leonard Niedrach, devised a way of depositing platinum onto the membrane, which served as catalyst for the necessary hydrogen oxidation and oxygen reduction reactions. This became known as the 'Grubb-Niedrach fuel cell'. GE went on to develop this technology with NASA and McDonnell Aircraft, leading to its use during Project Gemini. This was the first commercial use of a fuel cell. It wasn't until 1959 that British engineer Francis Thomas Bacon successfully developed a 5 kW stationary fuel cell. In 1959, a team led by Harry Ihrig built a 15 kW fuel cell tractor for Allis-Chalmers which was demonstrated across the US at state fairs. This system used potassium hydroxide as the electrolyte and compressed hydrogen and oxygen as the reactants. Later in 1959, Bacon and his colleagues demonstrated a practical five-kilowatt unit capable of powering a welding machine. In the 1960s, Pratt and Whitney licensed Bacon's U.S. patents for use in the U.S. space program to supply electricity and drinking water (hydrogen and oxygen being readily available from the spacecraft tanks).
United Technology Corp.'s UTC Power subsidiary was the first company to manufacture and commercialize a large, stationary fuel cell system for use as a co-generation power plant in hospitals, universities and large office buildings. UTC Power continues to market this fuel cell as the PureCell 200, a 200 kW system. UTC Power continues to be the sole supplier of fuel cells to NASA for use in space vehicles, having supplied the Apollo missions, and currently the Space Shuttle program, and is developing fuel cells for automobiles, buses, and cell phone towers; the company has demonstrated the first fuel cell capable of starting under freezing conditions with its proton exchange membrane automotive fuel cell.

Types Of Fuel Cells

There are several different types of fuel cells, each using a different chemistry. Fuel cells are usually classified by their operating temperature and the type of electrolyte they use. Some types of fuel cells work well for use in stationary power generation plants. Others may be useful for small portable applications or for powering cars. The main types of fuel cells include:
Polymer exchange membrane fuel cell (PEMFC)
The Department of Energy (DOE) is focusing on the PEMFC as the most likely candidate for transportation applications. The PEMFC has a high power density and a relatively low operating temperature (ranging from 60 to 80 degrees Celsius, or 140 to 176 degrees Fahrenheit). The low operating temperature means that it doesn't take very long for the fuel cell to warm up and begin generating electricity. We?ll take a closer look at the PEMFC in the next section.
Solid oxide fuel cell (SOFC)
These fuel cells are best suited for large-scale stationary power generators that could provide electricity for factories or towns. This type of fuel cell operates at very high temperatures (between 700 and 1,000 degrees Celsius). This high temperature makes reliability a problem, because parts of the fuel cell can break down after cycling on and off repeatedly. However, solid oxide fuel cells are very stable when in continuous use. In fact, the SOFC has demonstrated the longest operating life of any fuel cell under certain operating conditions. The high temperature also has an advantage: the steam produced by the fuel cell can be channeled into turbines to generate more electricity. This process is called co-generation of heat and power (CHP) and it improves the overall efficiency of the system.
Alkaline fuel cell (AFC)
This is one of the oldest designs for fuel cells; the United States space program has used them since the 1960s. The AFC is very susceptible to contamination, so it requires pure hydrogen and oxygen. It is also very expensive, so this type of fuel cell is unlikely to be commercialized.
Molten-carbonate fuel cell (MCFC)
Like the SOFC, these fuel cells are also best suited for large stationary power generators. They operate at 600 degrees Celsius, so they can generate steam that can be used to generate more power. They have a lower operating temperature than solid oxide fuel cells, which means they don't need such exotic materials. This makes the design a little less expensive.
Phosphoric-acid fuel cell (PAFC)
The phosphoric-acid fuel cell has potential for use in small stationary power-generation systems. It operates at a higher temperature than polymer exchange membrane fuel cells, so it has a longer warm-up time. This makes it unsuitable for use in cars.
Direct-methanol fuel cell (DMFC)
Methanol fuel cells are comparable to a PEMFC in regards to operating temperature, but are not as efficient. Also, the DMFC requires a relatively large amount of platinum to act as a catalyst, which makes these fuel cells expensive.
In the following section, we will take a closer look at the kind of fuel cell the DOE plans to use to power future vehicles -- the PEMFC.

Effeciency Of Fuel Cell

The efficiency of a fuel cell is dependent on the amount of power drawn from it. Drawing more power means drawing more current, which increases the losses in the fuel cell. As a general rule, the more power (current) drawn, the lower the efficiency. Most losses manifest themselves as a voltage drop in the cell, so the efficiency of a cell is almost proportional to its voltage. For this reason, it is common to show graphs of voltage versus current (so-called polarization curves) for fuel cells. A typical cell running at 0.7 V has an efficiency of about 50%, meaning that 50% of the energy content of the hydrogen is converted into electrical energy; the remaining 50% will be converted into heat. (Depending on the fuel cell system design, some fuel might leave the system unreacted, constituting an additional loss.)
For a hydrogen cell operating at standard conditions with no reactant leaks, the efficiency is equal to the cell voltage divided by 1.48 V, based on the enthalpy, or heating value, of the reaction. For the same cell, the second law efficiency is equal to cell voltage divided by 1.23 V. (This voltage varies with fuel used, and quality and temperature of the cell.) The difference between these numbers represents the difference between the reaction's enthalpy and Gibbs free energy. This difference always appears as heat, along with any losses in electrical conversion efficiency.
Fuel cells do not operate on a thermal cycle. As such, they are not constrained, as combustion engines are, in the same way by thermodynamic limits, such as Carnot cycle efficiency. At times this is misrepresented by saying that fuel cells are exempt from the laws of thermodynamics, because most people think of thermodynamics in terms of combustion processes (enthalpy of formation). The laws of thermodynamics also hold for chemical processes (Gibbs free energy) like fuel cells, but the maximum theoretical efficiency is higher (83% efficient at 298K) than the Otto cycle thermal efficiency (60% for compression ratio of 10 and specific heat ratio of 1.4). Comparing limits imposed by thermodynamics is not a good predictor of practically achievable efficiencies. Also, if propulsion is the goal, electrical output of the fuel cell has to still be converted into mechanical power with the corresponding inefficiency. In reference to the exemption claim, the correct claim is that the "limitations imposed by the second law of thermodynamics on the operation of fuel cells are much less severe than the limitations imposed on conventional energy conversion systems". Consequently, they can have very high efficiencies in converting chemical energy to electrical energy, especially when they are operated at low power density, and using pure hydrogen and oxygen as reactants.
For a fuel cell operating on air (rather than bottled oxygen), losses due to the air supply system must also be taken into account. This refers to the pressurization of the air and dehumidifying it. This reduces the efficiency significantly and brings it near to that of a compression ignition engine. Furthermore fuel cell efficiency decreases as load increases.
The tank-to-wheel efficiency of a fuel cell vehicle is about 45% at low loads and shows average values of about 36% when a driving cycle like the NEDC (New European Driving Cycle) is used as test procedure. The comparable NEDC value for a Diesel vehicle is 22%.
It is also important to take losses due to fuel production, transportation, and storage into account. Fuel cell vehicles running on compressed hydrogen may have a power-plant-to-wheel efficiency of 22% if the hydrogen is stored as high-pressure gas, and 17% if it is stored as liquid hydrogen.
Fuel cells cannot store energy like a battery, but in some applications, such as stand-alone power plants based on discontinuous sources such as solar or wind power, they are combined with electrolyzers and storage systems to form an energy storage system. The overall efficiency (electricity to hydrogen and back to electricity) of such plants (known as round-trip efficiency) is between 30 and 50%, depending on conditions While a much cheaper lead-acid battery might return about 90%, the electrolyzer/fuel cell system can store indefinite quantities of hydrogen, and is therefore better suited for long-term storage.
Solid-oxide fuel cells produce exothermic heat from the recombination of the oxygen and hydrogen. The ceramic can run as hot as 800 degrees Celsius. This heat can be captured and used to heat water in a micro combined heat and power (m-CHP) application. When the heat is captured, total efficiency can reach 80-90%. CHP units are being developed today for the European home market.

Design Issues And Advancements

• Costs. In 2002, typical cells had a catalyst content of US$1000 per kilowatt of electric power output. In 2008 UTC Power has 400kw Fuel cells for $1,000,000 per 400kW installed costs. The goal is to reduce the cost in order to compete with current market technologies including gasoline internal combustion engines. Many companies are working on techniques to reduce cost in a variety of ways including reducing the amount of platinum needed in each individual cell. Ballard Power Systems have experiments with a catalyst enhanced with carbon silk which allows a 30% reduction (1 mg/cm² to 0.7 mg/cm²) in platinum usage without reduction in performance. Monash University, Melbourne uses PEDOT instead of platinum.
• The production costs of the PEM (proton exchange membrane). The Nafion membrane currently costs €400/m². In 2005 Ballard Power Systems announced that its fuel cells will use Solupor, a porous polyethylene film patented by DSM.
• Water and air management (in PEMFCs). In this type of fuel cell, the membrane must be hydrated, requiring water to be evaporated at precisely the same rate that it is produced. If water is evaporated too quickly, the membrane dries, resistance across it increases, and eventually it will crack, creating a gas "short circuit" where hydrogen and oxygen combine directly, generating heat that will damage the fuel cell. If the water is evaporated too slowly, the electrodes will flood, preventing the reactants from reaching the catalyst and stopping the reaction. Methods to manage water in cells are being developed like electroosmotic pumps focusing on flow control. Just as in a combustion engine, a steady ratio between the reactant and oxygen is necessary to keep the fuel cell operating efficiently.
• Temperature management. The same temperature must be maintained throughout the cell in order to prevent destruction of the cell through thermal loading. This is particularly challenging as the 2H2 + O2 -> 2H2O reaction is highly exothermic, so a large quantity of heat is generated within the fuel cell.
• Durability, service life, and special requirements for some type of cells. Stationary fuel cell applications typically require more than 40,000 hours of reliable operation at a temperature of -35 °C to 40 °C (-31 °F to 104 °F), while automotive fuel cells require a 5,000 hour lifespan (the equivalent of 150,000 miles) under extreme temperatures. Automotive engines must also be able to start reliably at -30 °C (-22 °F) and have a high power to volume ratio (typically 2.5 kW per liter).

Fuel cell applications

Type 212 submarine with fuel cell propulsion of the German Navy in dock
Fuel cells are very useful as power sources in remote locations, such as spacecraft, remote weather stations, large parks, rural locations, and in certain military applications. A fuel cell system running on hydrogen can be compact and lightweight, and have no major moving parts. Because fuel cells have no moving parts and do not involve combustion, in ideal conditions they can achieve up to 99.9999% reliability. This equates to around one minute of down time in a two year period.
A new application is micro combined heat and power, which is cogeneration for family homes, office buildings and factories. The stationary fuel cell application generates constant electric power (selling excess power back to the grid when it is not consumed), and at the same time produces hot air and water from the waste heat. A lower fuel-to-electricity conversion efficiency is tolerated (typically 15-20%), because most of the energy not converted into electricity is utilized as heat. Some heat is lost with the exhaust gas just as in a normal furnace, so the combined heat and power efficiency is still lower than 100%, typically around 80%. In terms of exergy however, the process is inefficient, and one could do better by maximizing the electricity generated and then using the electricity to drive a heat pump. Phosphoric-acid fuel cells (PAFC) comprise the largest segment of existing CHP products worldwide and can provide combined efficiencies close to 90% (35-50% electric + remainder as thermal) Molten-carbonate fuel cells have also been installed in these applications, and solid-oxide fuel cell prototypes exist.

The world's first certified Fuel Cell Boat (HYDRA), in Leipzig/Germany
Since electrolyzer systems do not store fuel in themselves, but rather rely on external storage units, they can be successfully applied in large-scale energy storage, rural areas being one example. In this application, batteries would have to be largely oversized to meet the storage demand, but fuel cells only need a larger storage unit (typically cheaper than an electrochemical device).
One such pilot program is operating on Stuart Island in Washington State. There the Stuart Island Energy Initiative has built a complete, closed-loop system: Solar panels power an electrolyzer which makes hydrogen. The hydrogen is stored in a 500 gallon tank at 200 PSI, and runs a ReliOn fuel cell to provide full electric back-up to the off-the-grid residence. The SIEI website gives extensive technical details.
The world's first Fuel Cell Boat HYDRA used an AFC system with 6.5 kW net output.
Suggested applications
• Base load power plants
• Electric and hybrid vehicles.
• Auxiliary power
• Off-grid power supply
• Notebook computers for applications where AC charging may not be available for weeks at a time.
• Portable charging docks for small electronics (e.g. a belt clip that charges your cell phone or PDA).
• Smartphones with high power consumption due to large displays and additional features like GPS might be equipped with micro fuel cells.

Toyota FCHV PEM FC fuel cell vehicle
The first public hydrogen refueling station was opened in Reykjavík, Iceland in April 2003. This station serves three buses built by DaimlerChrysler that are in service in the public transport net of Reykjavík. The station produces the hydrogen it needs by itself, with an electrolyzing unit (produced by Norsk Hydro), and does not need refilling: all that enters is electricity and water. Royal Dutch Shell is also a partner in the project. The station has no roof, in order to allow any leaked hydrogen to escape to the atmosphere.
The GM 1966 Electrovan was the automotive industry's first attempt at an automobile powered by a hydrogen fuel cell. The Electrovan, which weighed more than twice as much as a normal van, could travel up to 70mph for 30 seconds
The 2001 Chrysler Natrium used its own on-board hydrogen processor. It produces hydrogen for the fuel cell by reacting sodium borohydride fuel with Borax, both of which Chrysler claimed were naturally occurring in great quantity in the United States. The hydrogen produces electric power in the fuel cell for near-silent operation and a range of 300 miles without impinging on passenger space. Chrysler also developed vehicles which separated hydrogen from gasoline in the vehicle, the purpose being to reduce emissions without relying on a nonexistent hydrogen infrastructure and to avoid large storage tanks.
In 2003 President George Bush proposed what is called the Hydrogen Fuel Initiative (HFI), which was later implemented by legislation through the 2005 Energy Policy Act and the 2006 Advanced Energy Initiative. These aim at further developing hydrogen fuel cells and its infrastructure technologies with the ultimate goal to produce fuel cell vehicles that are both practical and cost-effective by 2020. Thus far the United States has contributed 1 billion dollars to this project.
In 2005 the British firm Intelligent Energy produced the first ever working hydrogen run motorcycle called the ENV (Emission Neutral Vehicle). The motorcycle holds enough fuel to run for four hours, and to travel 100 miles in an urban area, at a top speed of 50 miles per hour. It will cost around $6,000 Honda is also going to offer fuel-cell motorcycles

A hydrogen fuel cell public bus accelerating at traffic lights in Perth, Western Australia
There are numerous prototype or production cars and buses based on fuel cell technology being researched or manufactured. Research is ongoing at a variety of motor car manufacturers. Honda has announced the release of a hydrogen vehicle in 2008.
Type 212 submarines use fuel cells to remain submerged for weeks without the need to surface.
Boeing researchers and industry partners throughout Europe are planning to conduct experimental flight tests in 2007 of a manned airplane powered only by a fuel cell and lightweight batteries. The Fuel Cell Demonstrator Airplane research project was completed recently and thorough systems integration testing is now under way in preparation for upcoming ground and flight testing. The Boeing demonstrator uses a Proton Exchange Membrane (PEM) fuel cell/lithium-ion battery hybrid system to power an electric motor, which is coupled to a conventional propeller.
Fuel cell powered race vehicles, designed and built by university students from around the world, competed in the world's first hydrogen race series called the 2008 Formula Zero Championship, which began on August 22nd, 2008 in Rotterdam, the Netherlands. The next race is in South Carolina in March 2009.
Not all geographic markets are ready for SOFC powered m-CHP appliances. Currently, the regions that lead the race in Distributed Generation and deployment of fuel cell m-CHP units are the EU and Japan.

Hydrogen economy

Electrochemical extraction of energy from hydrogen via fuel cells is an especially clean method of meeting power requirements, but not an efficient one, due to the necessity of adding large amounts of energy to either water or hydrocarbon fuels in order to produce the hydrogen. Additionally, during the extraction of hydrogen from hydrocarbons, carbon monoxide is released. Although this gas is artificially converted into carbon dioxide, such a method of extracting hydrogen remains environmentally injurious. It must however be noted that regarding the concept of the hydrogen vehicle, burning/combustion of hydrogen in an internal combustion engine (IC/ICE) is often confused with the electrochemical process of generating electricity via fuel cells (FC) in which there is no combustion (though there is a small byproduct of heat in the reaction). Both processes require the establishment of a hydrogen economy before they may be considered commercially viable, and even then, the aforementioned energy costs make a hydrogen economy of questionable environmental value. Hydrogen combustion is similar to petroleum combustion, and like petroleum combustion, still results in nitrogen oxides as a by-product of the combustion, which lead to smog. Hydrogen combustion, like that of petroleum, is limited by the Carnot efficiency, but is completely different from the hydrogen fuel cell's chemical conversion process of hydrogen to electricity and water without combustion. Hydrogen fuel cells emit only water during use, while producing carbon dioxide emissions during the majority of hydrogen production, which comes from natural gas. Direct methane or natural gas conversion (whether IC or FC) also generate carbon dioxide emissions, but direct hydrocarbon conversion in high-temperature fuel cells produces lower carbon dioxide emissions than either combustion of the same fuel (due to the higher efficiency of the fuel cell process compared to combustion), and also lower carbon dioxide emissions than hydrogen fuel cells, which use methane less efficiently than high-temperature fuel cells by first converting it to high purity hydrogen by steam reforming. Although hydrogen can also be produced by electrolysis of water using renewable energy, at present less than 3% of hydrogen is produced in this way.
Hydrogen is an energy carrier, and not an energy source, because it is usually produced from other energy sources via petroleum combustion, wind power, or solar photovoltaic cells. Hydrogen may be produced from subsurface reservoirs of methane and natural gas by a combination of steam reforming with the water gas shift reaction, from coal by coal gasification, or from oil shale by oil shale gasification. low pressure electrolysis of water or high pressure electrolysis, which requires electricity, and high-temperature electrolysis/thermochemical production, which requires high temperatures (ideal the for expected Generation IV reactors), are two primary methods for the extraction of hydrogen from water.
As of 2006, 49.0% of the electricity produced in the UnitedStates comes from coal, 19.4% comes from nuclear, 20.0% comes from natural gas, 7.0% from hydroelectricity, 1.6% from petroleum and the remaining 3.1% mostly coming from geothermal, solar and biomass. When hydrogen is produced through electrolysis, the energy comes from these sources. Though the fuel cell itself will only emit heat and water as waste, pollution is often caused when generating the electricity required to produce the hydrogen that the fuel cell uses as its power source (for example, when coal, oil, or natural gas-generated electricity is used). This will be the case unless the hydrogen is produced using electricity generated by hydroelectric, geothermal, solar, wind or other clean power sources (which may or may not include nuclear power, depending on one's attitude to the nuclear waste byproducts); hydrogen is only as clean as the energy sources used to produce it. A holistic approach has to take into consideration the impacts of an extended hydrogen scenario, including the production, the use and the disposal of infrastructure and energy converters.
Nowadays low temperature fuel cell stacks proton exchange membrane fuel cell (PEMFC), direct methanol fuel cell (DMFC) and phosphoric acid fuel cell (PAFC) make extensive use of catalysts. Impurities create catalyst poisoning (reducing activity and efficiency), thus high hydrogen purity or higher catalyst densities are required. Limited reserves of platinum quicken the synthesis of an inorganic complex. Although platinum is seen by some as one of the major "showstoppers" to mass market fuel cell commercialization companies, most predictions of platinum running out and/or platinum prices soaring do not take into account effects of thrifting (reduction in catalyst loading) and recycling. Recent research at Brookhaven National Laboratory could lead to the replacement of platinum by a gold-palladium coating which may be less susceptible to poisoning and thereby improve fuel cell lifetime considerably. Current targets for a transport PEM fuel cells are 0.2 g/kW Pt – which is a factor of 5 decrease over current loadings – and recent comments from major original equipment manufacturers (OEMs) indicate that this is possible. Also it is fully anticipated that recycling of fuel cells components, including platinum, will kick-in. High-temperature fuel cells, including molten carbonate fuel cells (MCFC's) and solid oxide fuel cells (SOFC's), do not use platinum as catalysts, but instead use cheaper materials such as nickel and nickel oxide, which are considerably more abundant (for example, nickel is used in fairly large quantities in common stainless steel).

Research and development

August 2005: Georgia Institute of Technology researchers use triazole to raise the operating temperature of PEM fuel cells from below 100 °C to over 125 °C, claiming this will require less carbon-monoxide purification of the hydrogen fuel.
2006: Staxon introduced an inexpensive OEM fuel cell module for system integration. In 2006 Angstrom Power, a British Columbia based company, began commercial sales of portable devices using proprietary hydrogen fuel cell technology, trademarked as "micro hydrogen

amorphous semiconducter and thin film preparation

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:29:00 -0700

A Progress Report
on
AMORPHOUS SEMICONDUCTOR
AND
THIN FILM PREPARATION
Submitted to: DR. KAILASH SIR
DEPTT. OF PHYSICS

Lovely Professional University, Jalandhar

“Bachelor of Technology”
(2008-09)
Submitted By

Rahul Mahajan & Ashwani Soni

Btech(CSE) Ist Sem

Uni Roll No.R237A27

Acknowledgement

Gratitude cannot be seen or expressed. It can only be felt in heart and is beyond description. Often words are inadequate to serve as a model of expression of one’s feeling, specially the sense of indebtedness and gratitude to all those who help us in our duty.

It is of immense pleasure and profound privilege to express my gratitude and indebtedness along with sincere thanks to DR. KAILASH SIR, Faculty of Lovely University for providing me the opportunity to work for a project on “AMORPHOUS SEMICONDUCTOR AND THIN FILM PREPARATION “

In particular I would like to mention the efforts of DR.KAILASH SIR, Lecturer of LPU Jalandhar, without whose encouragement the project could not have been started . He helped me on the project as an advisor and offered his help when needed in every aspect of project.

I am beholden to my family and friends for their blessings and encouragement.

Always Obediently
Rahul Mahajan

Introduction
Amorphous materials are of interest because of their complexity and unique structural and electronic properties. While crystalline and amorphous silicon are widely used in the manufacture of semiconductor devices, the carbon analogues at first seem of limited value in that context. This is due to the lack of electronic levels in the band gap of diamond on the introduction of n or p doping. However this is not true of the dense form of diamond-like amorphous carbon as opposed to the lower density graphitic form of amorphous carbon. To investigate the properties of such structures a model is first required of pure amorphous carbon. It is this problem which will be discussed in this chapter.
Of fundamental interest is the microscopic origins of such properties ranging from mechanical and elastic characteristics to the electronic and optical properties. As seen in Chapter 3, and in many recent publications (for example [115, 116, 117, 118]), carbon can display many different bonding configurations with varying coordination number due to the ability to form both bonding such as in graphitic structures and bonding as in diamond, although it seems energetically unfavourable to distort the bonding angles. Silicon also forms an amorphous phase which has generally been modelled using 4-fold coordinated continuous random networks of bonded atoms[115, 119]. These have used empirical and tight-binding force models which agree well with experiment. It is possible that this type of model may not be complete in view of the unusual interstitial configuration found in the previous chapter which was associated with a low defect formation energy. This implies that such a bonding topology could easily be formed in amorphous silicon[120].
In chapters 3 and 4, calculations were performed on complex forms of silicon and carbon which are characterised by short range order but still retain long range crystalline order. In view of the difficulties associated with performing full theoretical calculations on amorphous structures they proved to be a useful insight into the physics of short range disorder. It therefore seems a natural conclusion to attempt a molecular dynamics calculation on the amorphous structures of silicon and carbon and examine their differences.
Previous experimental and theoretical studies of the microscopic structure of amorphous carbon[121, 122, 123, 124] show that it is dependent on the macroscopic density which in turn depends on the method in which the sample was made. The trend in structures is from graphitic-like structures embedded in a matrix of both two-fold and four-fold coordinated atoms at a low density of 2.20 to 2.69 g/cm (found from tight binding molecular dynamics[124]), to diamond-like amorphous carbon containing `defected' three fold sites at a high density of 3.35 g/cm [121, 125]. This change in density also changes the bonding properties considerably[126], where the ratio of is found be inversely proportional to the density of the amorphous carbon.
Studies of amorphous silicon seem to show a somewhat simpler behaviour[116, 119] where the microscopic structure consists of distorted tetrahedral units. Numerous hand built and computer models have been constructed[127, 128, 129]. In the relaxed continuous random network models various potentials and bond charge models have been used (for example, the Keating and Stillinger-Weber potentials) in order to minimise total energies. Other models have included `defect' atoms that are three fold coordinated which have been obtained from various molecular dynamics techniques by cooling from the melt[130].
One of the main methods used recently for obtaining better models of both amorphous carbon and silicon[131] is that of reverse Monte Carlo simulations[132]. This method involves fitting the structure factors for trial atomic configurations to experimental results (measured, for example, by neutron diffraction) by moving the atoms at random and accepting the move by a probability given by the difference in the new structure factor and the experimental measurements. Configurations are accepted under certain constraints, such as bond length, coordination number, etc. Some configurations are not accepted such as those containing three membered rings. This may be incorrect given the results presented in Chapter 6 and the results for carbon given below. Also calculations on very small cells of silicon atoms arranged randomly have been found to contain 3-fold rings which was found to be relatively stable under small atomic displacements[133].

Structural Details

Silicon
Two samples of amorphous silicon have been generated at slightly different densities using the same initial random configuration. The higher density sample will be referred to as system-I, the other being system-II. After relaxation, the final configurations are found to be rather different. Given in Figures 8.1 and 8.2 are the radial distribution functions, g(r), of each sample.

Figure 8.1: Radial distribution function of amorphous silicon calculated at a density of 2.6g/cm .

Figure 8.2: Radial distribution function of amorphous silicon calculated at a density of 2.3g/cm .
The radial distribution function of system-I is in excellent agreement to that of experiment[131, 135]. For comparison, an experimental radial distribution function for an amorphous silicon sample which has a density of 2.45g/cm is given in Figure 8.3.

Figure 8.3: Experimental radial distribution function of amorphous silicon as found from neutron diffraction measurements.
System-II differs slightly in that the second neighbour peak at 3.5-4Å is too low and slightly wider than experiment. Integration under this part of the curve (taken from 2.9Å to 4.3Å) however indicates that the average number of nearest neighbours per atom is similar in each case. The lower density of system-II allows for a wider spread in second neighbour distances than system-I.
In both cases the height and width of the first peak of the radial distribution functions are in agreement with experiment. This implies that the average coordination number is correct in both cases. To find the coordination number, the maximum length of a silicon bond must be known. Unfortunately, the first minimum in the radial distribution function does not go to zero showing that there is a continuous range of neighbour distances. For this reason a maximum bond length is chosen arbitrarily to be 2.55Å. This then defines all the bonds within the structure. Some structural details of systems I and II are summarised in Table .
The coordination number of each atom can now be calculated. It is found that most atoms are four-fold coordinated (70.3% for system-I and 81.3% for system-II) while only a few atoms are either 3 or 5 fold coordinated. In system-I a single two-fold coordinated site is found. Such a feature has not been included in reverse Monte Carlo studies of amorphous silicon. In figure 8.4 is a schematic diagram of some 2, 3, and 4 fold coordinated atoms found in system-I. The 4-fold coordinated site is typical of the atomic structure of most of the atoms in the sample. It consists of a distorted tetrahedral bonding arrangement with a more distant, but unbonded, 5 neighbour, somewhat similar to that found in the BC8 structure. In the case shown in Figure 8.4, the bond angles for each site are given in Table .

Figure 8.4: Schematic diagram showing a group of atoms in the amorphous silicon structure of various coordination numbers. The full lines show covalent bonds while the dashed lines indicate unbonded close neighbouring atoms. The bonded/unbonded nature of neighbouring atoms is determined by examining 3d charge density plots. The distances are shown in Å.
The average bond angle for each coordination number is also shown in Table . As expected, the mean angle for the 4-fold coordinated sites are approximately that of the perfect tetrahedral angle. It may have been expected that the mean bond angle at a 3-fold site to be similar to that of an graphitic-like region, but instead it is found that it is less than . This leads to the implication that 3-fold sites are tending to have p-like character (at an angle of ). The bonding topology of the 3-fold site resembles a triangular pyramid with a well defined bond to the three neighbours of a central atom. A non-bonding orbital is formed at the top of the pyramid indicating that the `defect' site still retains bonding characteristics.
A typical five fold site is shown in Figure 8.5.

Figure 8.5: A five-fold coordinated silicon atom. The solid lines show the five covalent bonds from the central atom. The bond lengths are given in Å.
On examination of the valence electron charge density it is found that there are no bonds formed between atoms more that 2.6Å distant (hence the choice of when calculating the coordination number). There exist several sites with 5 atoms much closer than this (there are no 6-fold sites found) and hence covalent bonding is expected to occur. On examination of the valence electron charge density around these sited, we find that the sites tend to be fully 3 or 4 fold coordinated where the remaining atoms are relatively close and form slightly weaker bonds. When a fifth atom is found in an otherwise tetrahedral configuration it tends to weaken the longer bonds further. Such a configuration is shown in Figure 8.5 where the two more distant atoms form much weaker bonds that the three closer ones. Also, it is found that usually one of the neighbours of a five-fold coordinated atom has only three neighbours. This suggests that the electrons for the fifth bond is `donated' from the undercoordinated atom. Note that this is similar to the five-fold configuration found in the Si-BC8 surface in Chapter 5.
The average coordination number for system-I is 4.03 while the slightly lower density structure of system-II is found to have a coordination number of 3.97. Most models of amorphous silicon, such as random networks, assume from experimental measurements that the structure is fully four-fold coordinated. The fact that both simulations have found a structure which is very close to those found in other calculations from an initial random packing indicates that a full annealing treatment may not be necessary. There have been previous ab initio calculations on amorphous silicon and germanium (for example, see [135, 130]) which have rapidly cooled the melt in order to form models of their amorphous structure which reduces the percentage of `wrongly' coordinated sites, but giving an average coordination number similar to that found here. This suggests that cooling from the melt followed by annealing may not be the most efficient method of obtaining a reasonable model of amorphous silicon since there results do not differ significantly from those given here.
The ring statistics are also given in Table for both systems. This can be compared to the BC8 and ST12 structures which contain a range of small ring sizes from 5 to 7 fold rings. There are no three fold rings (and therefore no three-centre bonding orbitals) found in either sample which is somewhat unexpected considering the low formation energy of the interstitial configuration found in Chapter 6.
Since the bond lengths found in amorphous silicon (and also in the BC8 and ST12 structures) are similar to that found in the diamond structure it is generally assumed that the energy associated with straining the bond angle away from the perfect tetrahedral value gives the main proportion of excess energy of the amorphous structure relative to that of diamond Si[130]. This distortion of angles away from is illustrated in Figures 8.6 and 8.7 which show the bond angle distribution functions for systems I and II respectively.

Figure 8.6: Bond angle distribution function of amorphous silicon system-I.

Figure 8.7: Bond angle distribution function of amorphous silicon system-II.
As can be seen, a relatively small change in the size of the unit cell of the 64 atom simulation (a change of 3.5%) makes a rather large change to both the radial and bond angle distribution functions despite the same random starting configurations. The main features however seem to remain similar in both cases. There is a very large spread in bond angles centred about the maximum of .

Figure 8.8: A comparison between the first five neighbour distances in BC8, ST12 and amorphous silicon. The points on the BC8 and ST12 graphs show the neighbour distances at several different compressions.
The average first five neighbour distances for the two amorphous silicon simulations are shown along side a similar plot silicon in the BC8 and ST12 structures. Also shown is a plot of these distances for a third simulation of amorphous silicon as a much reduced volume to emphasize this trend. The neighbour distances for BC8 and ST12 are shown over a wide range of pressures. Firstly, it should be noticed that the distance to the first four bonded neighbours remains relatively unchanged with respect to the (generally unbonded) fifth neighbour distance. It can be seen that this is increasingly true in the trend of BC8 ST12 amorphous as the structure becomes more disordered. In the similar plot for a highly compressed amorphous silicon simulation the trend in neighbour distances becomes linear. It should be noted, however, that this third simulation is done only to show this reduction in the trend of reducing the distance in the extreme case where experimental verification of this linear trend in neighbour distances may be infeasible.
Carbon
The amorphous carbon structure is generated by the same method used above for silicon where the same initial random configuration is used. This allows a direct comparison between the amorphous silicon and carbon structures. Figures 8.9 and 8.10 show the radial and angle distribution functions of the final atomic configuration of the sample.

Figure 8.9: Radial distribution function of amorphous carbon calculated at a density of 3.4g/cm .

Figure 8.10: Bond angle distribution function of amorphous carbon.
It should be noticed that, unlike both silicon simulations, the radial distribution function drops to zero after the first neighbour peak. This leads to an unambiguous method of locating bonded pairs of atoms: 1.85Å. Examination of the electronic charge density indicates that this definition of a bond is correct. Using this the coordination numbers and ring statistics can be found. These are summarised in Table 8.1.

Table 8.1: Structural data for the amorphous carbon simulation.
As can be seen from the coordination numbers there exists no atom to which a fifth nearest neighbour is bonded. On consideration of the case of BC8 carbon this is expected. It was found that carbon is unable to form highly distorted tetrahedral bonding, favouring instead multiple bonding to a single atom. In order for silicon to form a 5-fold coordinated atom it is necessary to form a wide range of bond angles (from about to was found in the two silicon samples for ). Although the chemistry of carbon allows it to form many bonding configurations, this one is unstable with respect to multiply covalent bonds.
Although the radial distribution function is a very useful quantity in determining averages for shells of neighbouring atoms it is not unambiguously related to the spatial distribution of the carbon atoms. The bond angle distribution function is also necessary to determine the types of bonding. The bond angle distribution function contains several interesting features. Other experimental and theoretical results indicate that amorphous carbon contains mainly four-fold coordinated bonded atoms. This is also evident here in the large peak at about . Averaging the bond angles subtended by all four-fold coordinated atoms gives an angle of . Also of note is the shoulder at indicating planar graphitic-like bonding is also present, although in a smaller amount. Averaging the bond angles of three-fold coordinated atoms gives which is slightly less than the expected for perfect -like bonding although the statistics are rather limited since only six sites are found. There is also an indication that the amorphous carbon may be forming some p-like bonding due to the peak in the distribution function at . A small peak also appears at . Such a small bond angle indicates the possibility of 3-fold rings exist in the sample. A ring counting calculation in fact confirms that there are two such 3-fold rings (Table 8.1).
There have been several studies on the structure of amorphous silicon and carbon using the method of reverse Monte Carlo simulations which fits trial atomic configurations to the experimental radial distribution function. This may not fully describe all the atomic bonding environments in view of the large number of possible bonding topologies evident in the bond angle distribution function. Unfortunately obtaining this three-body function experimentally proves to be extremely difficult.
The amorphous structure is found to contain only three and four-fold coordinated atoms. Example of their bonding topologies are illustrated in Figure 8.11.

Figure 8.11: Schematic diagram illustrating a typical bonding topology of a 3-fold and 4-fold coordinated carbon atom. The dashed lines indicate close, but unbonded neighbours. The 4-fold coordinated atom forms a slightly distorted tetrahedral structure. It is this type of structure which dominates the amorphous carbons structure at this density.
There is a large number of 4-fold coordinated sites (90.6%) formed from slightly distorted bonding. All of the remaining sites are found to be 3-fold coordinated, but are not all necessary bonded. Of particular interest is the 3-fold rings that are found in the structure (see Figure 8.12).

Figure 8.12: A 3-fold ring of carbon atoms found in the amorphous structure. On examination of the charge density it was found that the electronic structure within the ring is best described as a 3-centre orbital, rather than three simple covalent bonds. Also shown is a four fold ring of carbon atoms is formed from a ring of covalent bonds unlike the 3-fold ring.
Like the silicon interstitial configuration found in Chapter 6 which formed a three fold ring, the charge density for this configuration in amorphous carbon formed a three-centre bonding orbital. Such a feature will not be found on examination of a radial distribution function alone since the inter-atomic distances are close to the C-C bond length. To find the electronic band(s) which are associated with the three-centre orbital, the electronic charge density was constructed from each individual band. This unusual feature is found to be very stable with its eigenvalue lying 24eV below the highest occupied band. However, another localised bonding orbital was found be be associated with it whose eigenvalue showed that it occupied the most energetic band.

Electronic properties
This section will discuss the electronic structure of amorphous silicon and carbon found in the above simulations. There have been many calculations on the electronic density of states of amorphous group IV materials in recent years, each result varying from the others depending on the model used to obtain the atomic coordinates[115, 116]. Amorphous carbon is atypical of the group IV semiconductors because of the large number of different bonding types that it can form. The electronic structure is governed by the relative importance of three and four fold sites. A purely four fold coordinated model of amorphous carbon[136] predicts only bonding to occur which gives a large gap in the electronic density of states. The electronic structure predicted by this model is similar to a broadened diamond-carbon density of states. It is now clear that this is not the correct model for diamond-like amorphous carbon and later tight-binding calculations[122, 125, 137] have found states which close the gap and have been associated to 3-fold coordinated atoms exhibiting bonding. The total number of states in the ` gap' increase when orbitals are introduced into the simulation. However, some models[137] produced from the Tersoff potential have a significant density of states near the Fermi level. This is in contradiction to experimental and other ab initio calculations[117, 125] which show only a small density of states at the Fermi level.
The electronic structure of the amorphous carbon simulation performed here is shown in Figure 8.13.

Figure 8.13: Electronic density of states for diamond-like amorphous carbon. Each of the different bonding regions are labelled.
The method detailed in Chapter 3 for calculating band structures is used here, where diagonalisation of the Hamiltonian matrix consisting of 128 occupied bands and a further 64 unoccupied bands is performed. The part of the density of states corresponding to bonding is very similar to a broadened diamond-like electronic structure. Most of the states around the Fermi level are found to be -like in nature leaving no band gap. Therefore the optical properties of amorphous carbon will be dominated by the bonded sites. There are however, relatively few of these (less than 10% of the atoms in the sample are 3-fold coordinated). They are not found to be clustered together as some earlier models of amorphous carbon predicted[116]. Instead they are found either bonded to three 4-fold coordinated atoms leaving a single electron in a localised p-like orbital, or rather often to two 4-fold atoms and another 3-fold site. This structure is similar to a recent ab initio calculation on diamond-like amorphous carbon where 3-fold sites are found to group in pairs[138]. Due to the lack of clustering of sites, it follows that it is the intermediate range correlations of the sites which will have profound effects on the optical spectrum.
It is also rather interesting to note that the density of states of diamond-like amorphous carbon calculated here is remarkably similar to ab initio calculations on the less dense graphitic form of amorphous carbon[117, 121, 122].
Contrary to that of carbon, the electronic structure of amorphous silicon is found to be predominantly composed of -bonding orbitals. It is now well established that the effects of structural disorder on -bonded tetrahedral systems are governed almost entirely by short ranged correlations. It is this fact that makes the `complex crystal model' of amorphous silicon a good approximation, but fails to do so in carbon where the medium range correlations play an important role.
The electronic density of states for systems-I and II of amorphous silicon are given in Figures 8.14 and 8.15 respectively.

Figure 8.14: Electronic density of states for system-I of amorphous silicon.

Figure 8.15: Electronic density of states for system-I of amorphous silicon.
On comparison of the two density of states diagrams, they are found to be very similar. Thus, a relatively large change in density and structure does little to change the electronic nature of the samples. Both exhibit a zero density at the Fermi level which agrees well with models based on random tetrahedral networks, although they are not in agreement with other ab initio calculations of amorphous silicon[130] which show a non-zero density at the Fermi level. This indicates that there are probably a large number of structural defects in their sample. However, experiment[139] shows a gap does exist in the density of states, in agreement with our calculations. On reconstruction of the charge density from the occupied bands around the Fermi level, these states are found to be localised mainly on the 3-fold coordinated atoms and can be classified as dangling bonds. Other localised states are found at atoms which are five-coordinated. These eigenstates are distributed throughout the five highly strained bonds at each 5-fold site. All of these localised states are found between the Fermi level and the large peak in the density of states at -2.5eV.

THIN FILM PREPARATION
The preparation of thin films of high quality is prerequisite for the realization of modern devices and sensors. Thus, the investigation and control of the growth process of new materials as thin films suitable for sensors is of great interest. We are working on superconducting, semiconducting, metallic, magnetic and ferroelectric materials for fundamental properties and applications in devices.
To deposit thin films we employ three different methods: Sputtering, thermal evaporation and pulsed laserdeposition (>>PLD) with in-situ Reflection High Energy Electron Diffraction (>>RHEED). The advantage of sputtering is its scalability to large areas. Thermal evaporation provides a simple means to prepare metallic contacts to devices. With PLD it is possible to prepare thin films from new oxide materials with interesting electronic properties in high quality within short time. The in-situ RHEED-analysis allows to control the thin film growth with monolayer precision up to a total film thickness of 1 um. Further analysis is done by x-ray diffraction (XRD) and scanning atomic force microscopy (>>AFM).
All optimization experiments are planned and evaluated with support of statistical methods for the design of experiments (DOE). To machieve highly reproducible thin films of high quality we develop further automatization in our processes. Also, we attach our thin film deposition equipment to the cleanroom area of class 100 by vacuum locks.

organic compunds

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:27:00 -0700

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What is an Organic Compound?
When you drive up to the pump at some gas stations you are faced with a variety of choices.
You can buy "leaded" gas or different forms of "unleaded" gas that have different octane numbers. As you filled the tank, you might wonder, "What is 'leaded' gas, and why do they add lead to gas?" Or, "What would I get for my money if I bought premium gas, with a higher octane number?"
You then stop to buy drugs for a sore back that has been bothering you since you helped a friend move into a new apartment. Once again, you are faced with choices (see the figure below). You could buy aspirin, which has been used for almost a hundred years. Or Tylenol, which contains acetaminophen. Or a more modern pain-killer, such as ibuprofen. While you are deciding which drug to buy, you might wonder, "What is the difference between these drugs?," and even, "How do they work?"

You then drive to campus, where you sit in a "plastic" chair to eat a sandwich that has been wrapped in "plastic," without worrying about why one of these plastics is flexibile while the other is rigid. While you're eating, a friend stops by and starts to tease you about the effect of your diet on the level of cholesterol in your blood, which brings up the questions, "What is cholesterol?" and "Why do so many people worry about it?"
Answers to each of these questions fall within the realm of a field known as organic chemistry. For more than 200 years, chemists have divided materials into two categories. Those isolated from plants and animals were classified as organic, while those that trace back to minerals were inorganic. At one time, chemists believed that organic compounds were fundamentally different from those that were inorganic because organic compounds contained a vital force that was only found in living systems.
The first step in the decline of the vital force theory occurred in 1828, when Friederich Wohler synthesized urea from inorganic starting materials. Wohler was trying to make ammonium cyanate (NH4OCN) from silver cyanate (AgOCN) and ammonium chloride (NH4Cl). What he expected is described by the following equation.
AgOCN(aq) + NH4Cl(aq) AgCl(s) + NH4OCN(aq)
The product he isolated from this reaction had none of the properties of cyanate compounds. It was a white, crystalline material that was identical to urea, H2NCONH2, which could be isolated from urine.

Neither Wohler nor his contemporaries claimed that his results disproved the vital force theory. But his results set in motion a series of experiments that led to the synthesis of a variety of organic compounds from inorganic starting materials. This inevitably led to the disappearance of "vital force" from the list of theories that had any relevance to chemistry, although it did not lead to the death of the theory, which still had proponents more than 90 years later.
If the difference between organic and inorganic compounds isn't the presence of some mysterious vital force required for their synthesis, what is the basis for distinguishing between these classes of compounds? Most compounds extracted from living organisms contain carbon. It is therefore tempting to identify organic chemistry as the chemistry of carbon. But this definition would include compounds such as calcium carbonate (CaCO3), as well as the elemental forms of carbon diamond and graphite that are clearly inorganic. We will therefore define organic chemistry as the chemistry of compounds that contain both carbon and hydrogen.
Even though organic chemistry focuses on compounds that contain carbon and hydrogen, more than 95% of the compounds that have isolated from natural sources or synthesized in the laboratory are organic. The special role of carbon in the chemistry of the elements is the result of a combination of factors, including the number of valence electrons on a neutral carbon atom, the electronegativity of carbon, and the atomic radius of carbon atoms (see the table below).
The Physical Properties of Carbon
Electronic configuration 1s2 2s2 2p2
Electronegativity 2.55
Covalent radius 0.077 nm
Carbon has four valence electrons 2s2 2p2 and it must either gain four electrons or lose four electrons to reach a rare-gas configuration. The electronegativity of carbon is too small for carbon to gain electrons from most elements to form C4- ions, and too large for carbon to lose electrons to form C4+ ions. Carbon therefore forms covalent bonds with a large number of other elements, including the hydrogen, nitrogen, oxygen, phosphorus, and sulfur found in living systems.
Because they are relatively small, carbon atoms can come close enough together to form strong C=C double bonds or even C C triple bonds. Carbon also forms strong double and triple bonds to nitrogen and oxygen. It can even form double bonds to elements such as phosphorus or sulfur that do not form double bonds to themselves.
Several years ago, the unmanned Viking spacecraft carried out experiments designed to search for evidence of life on Mars. These experiments were based on the assumption that living systems contain carbon, and the absence of any evidence for carbon-based life on that planet was presumed to mean that no life existed. Several factors make carbon essential to life.
• The ease with which carbon atoms form bonds to other carbon atoms.
• The strength of C C single bonds and the covalent bonds carbon forms to other nonmetals, such as N, O, P, and S.
• The ability of carbon to form multiple bonds to other nonmetals, including C, N, O, P, and S atoms.
These factors provide an almost infinite variety of potential structures for organic compounds, such as vitamin C shown in the figure below.

No other element can provide the variety of combinations and permutations necessary for life to exist.

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The Saturated Hydrocarbons or Alkanes
Compounds that contain only carbon and hydrogen are known as hydrocarbons. Those that contain as many hydrogen atoms as possible are said to be saturated. The saturated hydrocarbons are also known as alkanes.
The simplest alkane is methane: CH4. The Lewis structure of methane can be generated by combining the four electrons in the valence shell of a neutral carbon atom with four hydrogen atoms to form a compound in which the carbon atom shares a total of eight valence electrons with the four hydrogen atoms.

Methane is an example of a general rule that carbon is tetravalent; it forms a total of four bonds in almost all of its compounds. To minimize the repulsion between pairs of electrons in the four C H bonds, the geometry around the carbon atom is tetrahedral, as shown in the figure below.

Use the fact that carbon is usually tetravalent to predict the formula of ethane, the alkane that contains two carbon atoms.

The alkane that contains three carbon atoms is known as propane, which has the formula C3H8 and the following skeleton structure.

The four-carbon alkane is butane, with the formula C4H10.

The names, formulas, and physical properties for a variety of alkanes with the generic formula CnH2n+2 are given in the table below. The boiling points of the alkanes gradually increase with the molecular weight of these compounds. At room temperature, the lighter alkanes are gases; the midweight alkanes are liquids; and the heavier alkanes are solids, or tars.
The Saturated Hydrocarbons or Alkanes
Name Molecular
Formula Melting
Point (oC) Boiling
Point (oC) State
at 25oC
methane CH4 -182.5 -164 gas
ethane C2H6 -183.3 -88.6 gas
propane C3H8 -189.7 -42.1 gas
butane C4H10 -138.4 -0.5 gas
pentane C5H12 -129.7 36.1 liquid
hexane C6H14 -95 68.9 liquid
heptane C7H16 -90.6 98.4 liquid
octane C8H18 -56.8 124.7 liquid
nonane C9H20 -51 150.8 liquid
decane C10H22 -29.7 174.1 liquid
undecane C11H24 -24.6 195.9 liquid
dodecane C12H26 -9.6 216.3 liquid
excusing C20H42 36.8 343 solid
triacontane C30H62 65.8 449.7 solid
The alkanes in the table above are all straight-chain hydrocarbons, in which the carbon atoms form a chain that runs from one end of the molecule to the other. The generic formula for these compounds can be understood by assuming that they contain chains of CH2 groups with an additional hydrogen atom capping either end of the chain. Thus, for every n carbon atoms there must be 2n + 2 hydrogen atoms: CnH2n+2.
Because two points define a line, the carbon skeleton of the ethane molecule is linear, as shown in the figure below.

Because the bond angle in a tetrahedron is 109.5, alkanes molecules that contain three or four carbon atoms can no longer be thought of as "linear," as shown in the figure below.

Propane Butane
In addition to the straight-chain examples considered so far, alkanes also form branched structures. The smallest hydrocarbon in which a branch can occur has four carbon atoms. This compound has the same formula as butane (C4H10), but a different structure. Compounds with the same formula and different structures are known as isomers (from the Greek isos, "equal," and meros, "parts"). When it was first discovered, the branched isomer with the formula C4H10 was therefore given the name isobutane.
Isobutane
The best way to understand the difference between the structures of butane and isobutane is to compare the ball-and-stick models of these compounds shown in the figure below.

Butane Isobutane
Butane and isobutane are called constitutional isomers because they literally differ in their constitution. One contains two CH3 groups and two CH2 groups; the other contains three CH3 groups and one CH group.
There are three constitutional isomers of pentane, C5H12. The first is "normal" pentane, or n-pentane.

A branched isomer is also possible, which was originally named isopentane. When a more highly branched isomer was discovered, it was named neopentane (the new isomer of pentane).

Ball-and-stick models of the three isomers of pentane are shown in the figure below.

n-Pentane Isopentane

Neopentane

The following structures all have the same molecular formula: C6H14. Which of these structures represent the same molecule?

Determine the number of constitutional isomers of hexane, C6H14.

There are two constitutional isomers with the formula C4H10, three isomers of C5H12, and five isomers of C6H14. The number of isomers of a compound increases rapidly with additional carbon atoms. There are over 4 billion isomers for C30H62, for example.

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The Cycloalkanes
If the carbon chain that forms the backbone of a straight-chain hydrocarbon is long enough, we can envision the two ends coming together to form a cycloalkane. One hydrogen atom has to be removed from each end of the hydrocarbon chain to form the C C bond that closes the ring. Cycloalkanes therefore have two less hydrogen atoms than the parent alkane and a generic formula of CnH2n.
The smallest alkane that can form a ring is cyclopropane, C3H6, in which the three carbon atoms lie in the same plane. The angle between adjacent C C bonds is only 60, which is very much smaller than the 109.5 angle in a tetrahedron, as shown in the figure below.

Cyclopropane is therefore susceptible to chemical reactions that can open up the three-membered ring.
Any attempt to force the four carbons that form a cyclobutane ring into a plane of atoms would produce the structure shown in the figure below, in which the angle between adjacent C C bonds would be 90.

One of the four carbon atoms in the cyclobutane ring is therefore displaced from the plane of the other three to form a "puckered" structure that is vaguely reminiscent of the wings of a butterfly.
The angle between adjacent C C bonds in a planar cyclopentane molecule would be 108, which is close to the ideal angle around a tetrahedral carbon atom. Cyclopentane is not a planar molecule, as shown in the figure below, because displacing two of the carbon atoms from the plane of the other three produces a puckered structure that relieves some of the repulsion between the hydrogen atoms on adjacent carbon atoms in the ring.

By the time we get to the six-membered ring in cyclohexane, a puckered structure can be formed by displacing a pair of carbon atoms at either end of the ring from the plane of the other four members of the ring. One of these carbon atoms is tilted up, out of the ring, whereas the other is tilted down to form the "chair" structure shown in the figure below.

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Rotation Around C C Bonds
As one looks at the structure of the ethane molecule, it is easy to fall into the trap of thinking about this molecule as if it was static. Nothing could be further from the truth. At room temperature, the average velocity of an ethane molecule is about 500 m/s more than twice the speed of a Boeing 747. While it moves through space, the molecule is tumbling around its center of gravity like an airplane out of control. At the same time, the C H and C C bonds are vibrating like a spring at rates as fast as 9 x 1013 s-1.
There is another way in which the ethane molecule can move. The CH3 groups at either end of the molecule can rotate with respect to each around the C C bond. When this happens, the molecule passes through an infinite number of conformations that have slightly different energies. The highest energy conformation corresponds to a structure in which the hydrogen atoms are "eclipsed." If we view the molecule along the C C bond, the hydrogen atoms on one CH3 group would obscure those on the other, as shown in the figure below.

The lowest energy conformation is a structure in which the hydrogen atoms are "staggered," as shown in the figure below.

The difference between the eclipsed and staggered conformations of ethane is best illustrated by viewing these molecules along the C C bond, as shown in the figure below.

Eclipsed Staggered
The difference between the energies of these conformations is relatively small, only about 12 kJ/mol. But it is large enough that rotation around the C C bond is not smooth. Although the frequency of this rotation is on the order of 1010 revolutions per second, the ethane molecule spends a slightly larger percentage of the time in the staggered conformation.
The different conformations of a molecule are often described in terms of Newman projections. These line drawings show the six substituents on the C C bond as if the structure of the molecule was projected onto a piece of paper by shining a bright light along the C C bond in a ball-and-stick model of the molecule. Newman projections for the different staggered conformations of butane are shown in the figure below.

Because of the ease of rotation around C C bonds, there are several conformations of some of the cycloalkanes described in the previous section. Cyclohexane, for example, forms both the "chair" and "boat" conformations shown in the figure below.

Chair Boat
The difference between the energies of the chair conformation, in which the hydrogen atoms are staggered, and the boat conformation, in which they are eclipsed, is about 30 kJ/mol. As a result, even though the rate at which these two conformations interchange is about 1 x 105 s-1, we can assume that most cyclohexane molecules at any moment in time are in the chair conformation.

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The Nomenclature of Alkanes
Common names such as pentane, isopentane, and neopentane are sufficient to differentiate between the three isomers with the formula C5H12. They become less useful, however, as the size of the hydrocarbon chain increases.
The International Union of Pure and Applied Chemistry (IUPAC) has developed a systematic approach to naming alkanes and cycloalkanes based on the following steps.
• Find the longest continuous chain of carbon atoms in the skeleton structure. Name the compound as a derivative of the alkane with this number of carbon atoms. The following compound, for example, is a derivative of pentane because the longest chain contains five carbon atoms.

• Name the substituents on the chain. Substituents derived from alkanes are named by replacing the -ane ending with -yl. This compound contains a methyl (CH3-) substituent.

• Number the chain starting at the end nearest the first substituent and specify the carbon atoms on which the substituents are located. Use the lowest possible numbers. This compound, for example, is 2-methylpentane, not 4-methylpentane.

• Use the prefixes di-, tri-, and tetra- to describe substituents that are found two, three, or four times on the same chain of carbon atoms.
• Arrange the names of the substituents in alphabetical order.
Name the following compound.

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The Unsaturated Hydrocarbons: Alkenes and Alkynes
Carbon not only forms the strong C C single bonds found in alkanes, it also forms strong C=C double bonds. Compounds that contain C=C double bonds were once known as olefins (literally, "to make an oil") because they were hard to crystallize. (They tend to remain oily liquids when cooled.) These compounds are now called alkenes. The simplest alkenes have the formula C2H4 and the following Lewis structure.

The relationship between alkanes and alkenes can be understood by thinking about the following hypothetical reaction. We start by breaking the bond in an H2 molecule so that one of the electrons ends up on each of hydrogen atoms. We do the same thing to one of the bonds between the carbon atoms in an alkene. We then allow the unpaired electron on each hydrogen atom to interact with the unpaired electron on a carbon atom to form a new C H bond.

Thus, in theory, we can transform an alkene into the parent alkane by adding an H2 molecule across a C=C double bond. In practice, this reaction only occurs at high pressures in the presence of a suitable catalyst, such as piece of nickel metal.

Because an alkene can be thought of as a derivative of an alkane from which an H2 molecule has been removed, the generic formula for an alkene with one C=C double bond is CnH2n.
Alkenes are examples of unsaturated hydrocarbons because they have fewer hydrogen atoms than the corresponding alkanes. They were once named by adding the suffix -ene to the name of the substituent that carried the same number of carbon atoms.

The IUPAC nomenclature for alkenes names these compounds as derivatives of the parent alkanes. The presence of the C=C double bond is indicated by changing the -ane ending on the name of the parent alkane to -ene.

The location of the C=C double bond in the skeleton structure of the compound is indicated by specifying the number of the carbon atom at which the C=C bond starts.

The names of substituents are then added as prefixes to the name of the alkene.

Compounds that contain C C triple bonds are called alkynes. These compounds have four less hydrogen atoms than the parent alkanes, so the generic formula for an alkyne with a single C C triple bond is CnH2n-2. The simplest alkyne has the formula C2H2 and is known by the common name acetylene.

The IUPAC nomenclature for alkynes names these compounds as derivatives of the parent alkane, with the ending -yne replacing -ane.

In addition to compounds that contain one double bond (alkenes) or one triple bond (alkynes), we can also envision compounds with two double bonds (dienes), three double bonds (trienes), or a combination of double and triple bonds.

Chemical structure of methane, the simplest alkane
Alkanes, also known as paraffins are chemical compounds that consist only of the elements carbon (C) and hydrogen (H) (i.e., hydrocarbons), wherein these atoms are linked together exclusively by single bonds (i.e., they are saturated compounds) without any cyclic structure (i.e. loops). Alkanes belong to a homologous series of organic compounds in which the members differ by a constant relative atomic mass of 14.
Each carbon atom must have 4 bonds (either C-H or C-C bonds), and each hydrogen atom must be joined to a carbon atom (H-C bonds). A series of linked carbon atoms is known as the carbon skeleton or carbon backbone. In general, the number of carbon atoms is often used to define the size of the alkane (e.g., C2-alkane).
An alkyl group is a functional group or side-chain that, like an alkane, consists solely of singly-bonded carbon and hydrogen atoms, for example a methyl or ethyl group.
Saturated hydrocarbons can be linear (general formula CnH2n+2) wherein the carbon atoms are joined in a snake-like structure, branched (general formula CnH2n+2, n>3) wherein the carbon backbone splits off in one or more directions, or cyclic (general formula CnH2n, n>2) wherein the carbon backbone is linked so as to form a loop. According to the definition by IUPAC, the former two are alkanes, whereas the third group is called cycloalkanes. In other words, saturated hydrocarbons are divided into alkanes and cycloalkanes, depending on whether or not they have cyclic structures, and, in the technical sense, cycloalkanes are not alkanes. However, cycloalkanes are sometimes called cyclic alkanes, which can be confusing when "real" alkanes are called acyclic alkanes. Saturated hydrocarbons can also combine any of the linear, cyclic (e.g., polycyclic) and branching structures, and they are still alkanes (no general formula) as long as they are acyclic (i.e., having no loops).
The simplest possible alkane (the parent molecule) is methane, CH4. There is no limit to the number of carbon atoms that can be linked together, the only limitation being that the molecule is acyclic, is saturated, and is a hydrocarbon. Saturated oils and waxes are examples of larger alkanes where the number of carbons in the carbon backbone tends to be greater than 10.
Alkanes are not very reactive and have little biological activity. Alkanes can be viewed as a molecular scaffold upon which can be hung the interesting biologically-active/reactive portions (functional groups) of the molecule.

Different C4-alkanes and -cycloalkanes (left to right) n-butane and isobutane are the two C4H10 isomers; cyclobutane and methylcyclopropane are the two C4H8 isomers; bicyclo[1.1.0]butane is the only C4H6 isomer; tetrahedrane (not shown) is the only C4H4 isomer.
Alkanes with more than three carbon atoms can be arranged in a multiple number of ways, forming different structural isomers. An isomer is like a chemical anagram, in which the atoms of a chemical compound are arranged or joined together in a different order. The simplest isomer of an alkane is the one in which the carbon atoms are arranged in a single chain with no branches. This isomer is sometimes called the n-isomer (n for "normal", although it is not necessarily the most common). However the chain of carbon atoms may also be branched at one or more points. The number of possible isomers increases rapidly with the number of carbon atoms (sequence A000602 in OEIS). For example:
• C1: 1 isomer—methane
• C2: 1 isomer—ethane
• C3: 1 isomer—propane
• C4: 2 isomers—, n-butane isobutane
• C12: 355 isomers
• C32: 27,711,253,769 isomers
• C60: 22,158,734,535,770,411,074,184 isomers, many of which are not stable.
Branched alkanes can be chiral: 3-Methylhexane and its higher homologues are chiral due to their stereogenic center at carbon atom number 3. Chiral alkanes are of certain importance in biochemistry, as they occur as sidechains in chlorophyll and tocopherol (vitamin E). Chiral alkanes can be resolved into their enantiomers by enantioselective chromatography.
In addition to these isomers, the chain of carbon atoms may form one or more loops. Such compounds are called cycloalkanes.
Nomenclature
Main article: Organic nomenclature
The IUPAC nomenclature (systematic way of naming compounds) for alkanes is based on identifying hydrocarbon chains. Unbranched, saturated hydrocarbon chains are named systematically with a Greek numerical prefix denoting the number of carbons and the suffix "-ane".August Wilhelm von Hofmann suggested systematizing nomenclature by using the whole sequence of vowels a, e, i, o and u to create suffixes -ane, -ene, -ine (or -yne), -one, -une, for the hydrocarbons. The first three name hydrocarbons with single, double and triple bonds; "-one" represents a ketone; "-ol" represents an alcohol or OH group; "-oxy-" means an ether and refers to oxygen between two carbons, so that methoxy-methane is the IUPAC name for dimethyl ether.
It is difficult or impossible to find compounds with more than one IUPAC name. This is because shorter chains attached to longer chains are prefixes and the convention includes brackets. Numbers in the name, referring to which carbon a group is attached to, should be as low as possible, so that 1- is implied and usually omitted from names of organic compounds with only one side-group; "1-" is implied in Nitro-octane. Symmetric compou will have two ways of arriving at the same name.
Straight-chain alkanes are sometimes indicated by the prefix n- (for normal) where a non-linear isomer exists. Although this is not strictly necessary, the usage is still common in cases where there is an important difference in properties between the straight-chain and branched-chain isomers, e.g., n-hexane or 2- or 3-methylpentane.
The first four members of the series (in terms of number of carbon atoms) are named as follows:
methane, CH4
ethane, C2H6
propane, C3H8
butane, C4H10
Alkanes with five or more carbon atoms are named by adding the suffix -ane to the appropriate Greek-language prefix numerical multiplier with elision of any terminal vowel (-a or -o) from the basic numerical term. Hence, pentane, C5H12; hexane, C6H14; heptane, C7H16; octane, C8H18; etc. For a more complete list, see List of alkanes.
Branched alkanes

Ball-and-stick model of isopentane (common name) or 2-methylbutane (IUPAC systematic name)
Simple branched alkanes often have a common name using a prefix to distinguish them from linear alkanes, for example n-pentane, isopentane, and neopentane.
IUPAC naming conventions can be used to produce a systematic name.
The key steps in the naming of more complicated branched alkanes are as follows: Identify the longest continuous chain of carbon atoms
• Name this longest root chain using standard naming rules
• Name each side chain by changing the suffix of the name of the alkane from "-ane" to "-yl"
• Number the root chain so that sum of the numbers assigned to each side group will be as low as possible
• Number and name the side chains before the name of the root chain
• If there are multiple side chains of the same type, use prefixes such as "di-" and "tri-" to indicate it as such, and number each one.
Comparison of nomenclatures for three isomers of C5H12
Common name n-pentane isopentane neopentane
IUPAC name pentane 2-methylbutane 2,2-dimethylpropane
Structure

Main article: Cycloalkane
So-called cyclic alkanes are, in the technical sense, not alkanes, but cycloalkanes. They are hydrocarbons just like alkanes, but contain one or more rings.
Simple cycloalkanes have a prefix "cyclo-" to distinguish them from alkanes. Cycloalkanes are named as per their acyclic counterparts with respect to the number of carbon atoms, e.g., cyclopentane (C5H10) is a cycloalkane with 5 carbon atoms just like pentane (C5H12), but they are joined up in a five-membered ring. In a similar manner, propane and cyclopropane, butane and cyclobutane, etc.
Substituted cycloalkanes are named similar to substituted alkanes — the cycloalkane ring is stated, and the substituents are according to their position on the ring, with the numbering decided by Cahn-Ingold-Prelog rules.
Trivial names
The trivial (non-systematic) name for alkanes is "paraffins." Together, alkanes are known as the paraffin series. Trivial names for compounds are usually historical artifacts. They were coined before the development of systematic names, and have been retained due to familiar usage in industry. Cycloalkanes are also called naphthenes.
It is almost certain that the term paraffin stems from the petrochemical industry. Branched-chain alkanes are called isoparaffins . The use of the term "paraffin" is a general term and often does not distinguish between a pure compounds and mixtures of isomers with the same chemical formula (i.e., like a chemical anagram), e.g., pentane and isopentane.
Examples
The following trivial names are retained in the IUPAC system:
• isobutane for 2-methylpropane
• isopentane for 2-methylbutane
• neopentane for 2,2-dimethylpropane
Occurrence
Occurrence of alkanes in the Universe

Methane and ethane make up a large proportion of Jupiter's atmosphere

Extraction of oil, which contains many different hydrocarbons including alkanes
Alkanes form a significant portion of the atmospheres of the outer gas planets such as Jupiter (0.1% methane, 0.0002% ethane), Saturn (0.2% methane, 0.0005% ethane), Uranus (1.99% methane, 0.00025% ethane) and Neptune (1.5% methane, 1.5 ppm ethane). Titan (1.6% methane), a satellite of Saturn, was examined by the Huygens probe, which indicate that Titan's atmosphere periodically rains liquid methane onto the moon's surface. Also on Titan, a methane-spewing volcano was spotted and this volcanism is believed to be a significant source of the methane in the atmosphere. There also appear to be Methane/Ethane lakes near the north polar regions of Titan, as discovered by Cassini's radar imaging. Methane and ethane have also been detected in the tail of the comet Hyakutake. Chemical analysis showed that the abundances of ethane and methane were roughly equal, which is thought to imply that its ices formed in interstellar space, away from the Sun, which would have evaporated these volatile molecules. Alkanes have also been detected in meteorites such as carbonaceous chondrites.
Occurrence of alkanes on Earth
Traces of methane gas (about 0.0001% or 1 ppm) occur in the Earth's atmosphere, produced primarily by organisms such as Archaea, found for example in the gut of cows.
The most important commercial sources for alkanes are natural gas and oil. Natural gas contains primarily methane and ethane, with some propane and butane: oil is a mixture of liquid alkanes and other hydrocarbons. These hydrocarbons were formed when dead marine animals and plants (zooplankton and phytoplankton) died and sank to the bottom of ancient seas and were covered with sediments in an anoxic environment and converted over many millions of years at high temperatures and high pressure to their current form. Natural gas resulted thereby for example from the following reaction:
C6H12O6 → 3CH4 + 3CO2
These hydrocarbons collected in porous rocks, located beneath an impermeable cap rock and so are trapped. Unlike methane, which is constantly reformed in large quantities, higher alkanes (alkanes with 9 or more carbon atoms) rarely develop to a considerable extent in nature. These deposits, e.g., oil fields, have formed over millions of years and once exhausted cannot be readily replaced. The depletion of these hydrocarbons is the basis for what is known as the energy crisis.
Solid alkanes are known as tars and are formed when more volatile alkanes such as gases and oil evaporate from hydrocarbon deposits. One of the largest natural deposits of solid alkanes is in the asphalt lake known as the Pitch Lake in Trinidad and Tobago.
Methane is also present in what is called biogas, produced by animals and decaying matter, which is a possible renewable energy source.
Alkanes have a low solubility in water, so the content in the oceans is negligible; however, at high pressures and low temperatures (such as at the bottom of the oceans), methane can co-crystallize with water to form a solid methane hydrate.Although this cannot be commercially exploited at the present time, the amount of combustible energy of the known methane hydrate fields exceeds the energy content of all the natural gas and oil deposits put together;methane extracted from methane hydrate is considered therefore a candidate for future fuels.

Basic theory of semiconducters

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:25:00 -0700

SUPERCONDUCTORS

INDEX

 INTRODUCTION

 REVIEW OF LITERATURE

 IDEAS LEADING TO THE BCS THEORY

 BCS THEORY OF SUPERCONDUCTIVITY

 MATHEMATICAL NOTATION OF BCS THEORY

 SCIENTISTS WHO WORKED FOR BCS THEORY

 APPLICATIONS OF BCS THEORY

 FUTURE SCOPE

 BIBLIOGRAPHY

INTRODUCTION

I Ravneet kaur of B.tech IT(hons) department working on the topic called “ BCS Theory ” The topic of my term paper include mainly the approach towards the BCS theory of superconductors ie what is BCS theory, how it could be explained, Review of literature for BCS Theory and contribution of five scientists in working of BCS theory with applications and future scope of BCS Theory.

REVIEW OF LITERATURE
The BCS Theory of Superconductivity:
"a complete theoretical explanation of the phemonenon",

“The phenomenon of superconductivity was discovered by the Dutch physicist Kamerling Onnes in 1911. Already his first measurements indicated that one had found a fundamentally new state of matter... Many remarkable properties were discovered in the following decades. However, the central problem, the question about the underlying mechanism for superconductivity, remained a mystery up to the late 50's... A significant step forward was taken around 1950 when it was found theoretically and experimentally that the mechanism for superconductivity had to do with the coupling of electrons to the vibrations of the crystal lattice. Starting from this mechanism, Bardeen, Cooper and Schrieffer developed in 1957 a theory of superconductivity, which gave a complete theoretical explanation of the phenomenon.

[The BCS Theory] was indeed very successful in explaining in considerable detail the properties of superconductors. The theory also predicted new effects and it stimulated an intensive activity in theoretical and experimental research, which opened up new areas for research. One may as examples mention the use of the quantum mechanical tunnel phenomena to study superconductors, the discovery of magnetic flux quantization and the remarkable Josephson effects. These more recent developments are intimately connected with the fundamental theory of superconductivity and have confirmed in a striking way the validity of the theoretical concepts and ideas developed by Bardeen, Cooper and Schrieffer."

Ideas Leading to the BCS Theory
The BCS theory of superconductivity has successfully described the measured properties of Type I superconductors. It envisions resistance-free conduction of coupled pairs of electrons called Cooper pairs. This theory is remarkable enough that it is interesting to look at the chain of ideas which led to it.
1. One of the first steps toward a theory of superconductivity was the realization that there must be a band gap separating the charge carriers from the state of normal conduction.
o A band gap was implied by the very fact that the resistance is precisely zero. If charge carriers can move through a crystal lattice without interacting at all, it must be because their energies are quantized such that they do not have any available energy levels within reach of the energies of interaction with the lattice.
o A band gap is suggested by specific heats of materials like vanadium. The fact that there is an exponentially increasing specific heat as the temperature approaches the critical temperature from below implies that thermal
energy being used to bridge some kind of gap in energy. As the temperature increases, there is an exponential increase in the number of particles which would have enough energy to cross the gap.

2. The critical temperature for superconductivity must be a measure of the band gap, since the material could lose superconductivity if thermal energy could get charge carriers across the gap.

3. The critical temperature was found to depend upon isotopic mass. It certainly would not if the conduction was by free electrons alone. This made it evident that the superconducting transition involved some kind of interaction with the crystal lattice.

4. Single electrons could be eliminated as the charge carriers in superconductivity since with a system of fermions you don't get energy gaps. All available levels up to the Fermi energy fill up.

5. The needed boson behavior was consistent with having coupled pairs of electrons with opposite spins. The isotope effect described above suggested that the coupling mechanism involved the crystal lattice, so this gave rise to the phonon model of coupling envisioned with Cooper pairs.

BCS Theory of Superconductivity
An intuitive description of superconductivity is sufficient for public and non-technical uses. However, ultimately, a more rigorous mathematically- based explanation must be formulated. Superconductivity was not sufficiently explained until 1957 when John Bardeen and his graduate assistants Leon Cooper and John Schreiffer proposed a microscopic explanation that would later be their namesake: the BCS Theory. This theoretical explanation later earned them the Nobel prize, making John Bardeen the only man in history to be awarded this honor twice.

The BCS Theory is, in its simplest form, actually contradictory to our crude macroscopic view expressed earlier. As discussed earlier, superconductivity arises because electrons do not interact destructively with atoms in the crystal lattice of the material. The BCS Theory says that electrons do actually interact with the atoms, but constructively.

The BCS Theory makes a crucial assumption at the beginning: that an attractive force exists between electrons. In typical Type I superconductors, this force is due to Coulomb attraction between the electron and the crystal lattice. An electron in the lattice will cause a slight increase in positive charges around it. This increase in positive charge will, in turn, attract another electron. These two electrons are known as a Cooper pair. If the energy required to bind these electrons together is less than the energy from the thermal vibrations of the lattice attempting to break them apart, the pair will remain bound. This explains (roughly) why superconductivity requires low temperatures- the thermal vibration of the lattice must be small enough to allow the forming of Cooper pairs. In a superconductor, the current is made up of these Cooper pairs, rather than individual electrons.

So, Cooper pairs are formed by Coulomb interactions with the crystal lattice. This is also what overcomes resistance. Remember, an electron inside the lattice causes a slight increase of positive charge due to Coulomb attraction. As the Cooper pair flows, the leading electron causes this increase of charge, and the trailing electron is attracted by it. This is illustrated below.

This BCS theory prediction of Cooper pair interaction with the crystal lattice has been verified experimentally by the isotope effect. That is, the critical temperature of a material depends on the mass of the nucleus of the atoms. If an isotope is used (neutrons are added to make it more massive), the critical temperature decreases. This effect is most evident in Type I, and appears only weakly in Type II.
"...recall that early researchers made the somewhat paradoxical observation that the best conducting materials could not be made to exhibit superconductivity. A good conductor is, by definition, a material that will allow electrons to carry current with a minimum resistance. Therefore, since the primary cause of resistance is the electrons collisions with the lattice, a good conductor must have a minimal interaction between the electrons and the lattice. Consequently, the lattice is unable to mediate an attractive force between the electrons and the superconducting phase transition cannot occur. The converse of this observation also holds: metals exhibiting poor conductivity make excellent superconductors with relatively higher critical temperatures because the electrons greatly interact with the lattice." (Orlando 527)
This superconductivity of Cooper pairs is somewhat related to Bose-Einstein Condensation. The Cooper pairs act somewhat like bosons, which condense into their lowest energy level below the critical temperature, and lose electrical resistance.
The BCS Theory did exactly what a physical theory should do: it explained properties already witnessed in experiment, and it predicted experimentally verifiable phenomena. Though its specific quantitative elements were quite limited in their application (it only explained Type I s-wave superconductivity), its essence was quite broad and has been modified applied to various other superconductors, such as Type II perovskites.

Applications of BCS Theory
1. An induced superconducting state caused by charge transfer between intrinsically superconducting (α) and intrinsically normal (β) subsystems is studied. A most interesting case is a layered system with some layers being normal. An analysis of the general Hamiltonian describing the phenomenon allows us to evaluate Tc and the spectrum, which displays a two-gap structure. A superconducting state can be induced through different charge transfer channels (intrinsic proximity effect; inelastic two-band channel). A very important contribution comes from the ‘‘mixed’’ channel. Systems with various strengths of the coupling are described. The presence of magnetic impurities leads to an induced gapless superconductivity. This model is applied to high-Tc cuprates (in particular, to Y-Ba-Cu-O), as well as to conventional systems. The spectroscopy of Y-Ba-Cu-O appears to be very sensitive to the oxygen content whereas Tc changes relatively slowly. The model is directly related to such phenomena as residual microwave losses, zero-bias anomalies, the ‘‘plateau’’ effect, etc.
2. An induced superconducting state caused by charge transfer between intrinsically superconducting (α) and intrinsically normal (β) subsystems is studied. A most interesting case is a layered system with some layers being normal. An analysis of the general Hamiltonian describing the phenomenon allows us to evaluate Tc and the spectrum, which displays a two-gap structure. A superconducting state can be induced through different charge transfer channels (intrinsic proximity effect; inelastic two-band channel). A very important contribution comes from the ‘‘mixed’’ channel. Systems with various strengths of the coupling are described. The presence of magnetic impurities leads to an induced gapless superconductivity. This model is applied to high-Tc cuprates (in particular, to Y-Ba-Cu-O), as well as to conventional systems. The spectroscopy of Y-Ba-Cu-O appears to be very sensitive to the oxygen content whereas Tc changes relatively slowly. The model is directly related to such phenomena as residual microwave losses, zero-bias anomalies, the ‘‘plateau’’ effect, etc.

3. We propose a simple theory for the I-V curves of normal-superconducting microconstriction contacts which describes the crossover from metallic to tunnel junction behavior. The detailed calculations are performed within a generalized semiconductor model, with the use of the Bogoliubov equations to treat the transmission and reflection of particles at the N-S interface. By including a barrier of arbitrary strength at the interface, we have computed a family of I-V curves ranging from the tunnel junction to the metallic limit. Excess current, generated by Andreev reflection, is found to vary smoothly from 4Δ / 3eRN in the metallic case to zero for the tunnel junction. Charge-imbalance generation, previously calculated only for tunnel barriers, has been recalculated for an arbitrary barrier strength, and detailed insight into the conversion of normal current to supercurrent at the interface is obtained. We emphasize that the calculated differential conductance offers a particularly direct experimental test of the predictions of the model.

4. A new scheme to study exotic phase transitions is formulated by introducing the concept of a super-effective field. A general mechanism of phase transitions is elucidated and a general criterion of order parameters is proposed on the basis of the newly formulated super-effective-field theory. An alternative formulation based on a decoupled density matrix is also given. It is easily shown using these formulations that a quantum chiral order appears in the antiferromagnetic XY model on the triangular lattice. A super-effective-field theory of spin glasses is also presented. ©1988 The Physical Society of Japan

SCIENTISTS WHO WORKED FOR BCS THEORY

Dr. Bardeen's main fields of research since 1945 have been electrical conduction in semiconductors and metals, surface properties of semiconductors, theory of superconductivity, and diffusion of atoms in solids. The Nobel Prize in Physics was awarded in 1956 to John Bardeen, Walter H. Brattain, and William Shockley for "investigations on semiconductors and the discovery of the transistor effect," carried on at the Bell Telephone Laboratories. In 1957, Bardeen and two colleagues, L.N. Cooper and J.R. Schrieffer, proposed the first successful explanation of superconductivity, which has been a puzzle since its discovery in 1908. Much of his research effort since that time has been devoted to further extensions and applications of the theory.

Professor Cooper has received many forms of recognition for his work in 1972, he received the Nobel Prize in Physics (with J. Bardeen and J.R. Schrieffer) for his studies on the theory of superconductivity completed while still in his 20s. In 1968, he was awarded the Comstock Prize (with J.R. Schrieffer) of the National Academy of Sciences. The Award of Excellence, Graduate Faculties Alumni of Columbia University and Descartes Medal, Academie de Paris, Université Rene Descartes were conferred on Professor Cooper in the mid 1970s. In 1985, Professor Cooper received the John Jay Award of Columbia College. He holds seven honorary doctorates.

.
He served as Director of the Institute for Theoretical Physics in Santa Barbara from 1984-89. In 1992 he was appointed University Professor at Florida State University and Chief Scientist of the National High Magnetic Field Laboratory.He holds honorary degrees from the Technische Hochschule, Munich and the Universities of Geneva, Pennsylvania, Illinois, Cincinnati, Tel-Aviv, Alabama. In 1969 he was appointed by Cornell to a six-year term as a Andrew D. White Professor-at-Large.He is a member of the American Academy of Arts and Sciences, the National Academy of Sciences of which he is a member of their council, the American Philosophical Society, the Royal Danish Academy of Sciences and Letters and the Academy of Sciences of the USSR.The main thrust of his recent work has been in the area of high-temperature superconductivity, strongly correlated electrons, and the dynamics of electrons in strong magnetic fields.

FUTURE SCOPE

The ideas from which this area has quite recently emerged can be traced back at least some 40 years. Then, and quite independently, Larkin and Ovchinnikov (LO) and Fulde and Ferrell (FF) proposed on purely theoretical grounds what amounted to a new type of superconductivity, now often referred to as the LOFF phase. This phase of inhomogeneous superconductivity can be usefully viewed as a proposed generalization of the Bardeen-Cooper-Schrieffer (BCS) state which is appropriate to describe many properties
of elemental metallic superconductors. Whereas the basic building block of the BCS theory is the Cooper pair, where the two electrons have momenta equal in magnitude and opposite in direction,
in the so-called LOFF phase a salient feature is that momenta do not add to zero. Then an almost immediate consequence of the LOFF proposal
is that the energy gap, or order parameter, has a spatial variation.The LOFF proposal, to our knowledge, has not yet been confirmed beyond reasonable doubt in condensed matter, but expectations are high that such a phase will come up in real materials in the fore see able future. But what seems remarkable, and worthy of much fuller exploration, is that the same basic ideas of LOFF may also prove to play an important role in the future in nuclear physics and in the theory of some aspects of the properties of pulsars which are commonly identified with neutron stars.

BIBLIOGRAPHY

www.manhattanrarebooks-science.com
www.hyperphysics.com
www.ffden.com
www.wikipedia.com
www.brittanica.com
www.quench_analysis.com

blue ray

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:23:00 -0700

r
Topics Pages

1.Review of Litrature 3
2.Introduction 4
3.History & Origin 6
4.Blu-ray Association & Members 7
5.Cotributors 9
6.Launch and Sales Development 10
7.Competion from HD DVD 10
8.Technical Specification 11
9.Laser Optics 12
10.Hard-Coating Technology 12
11.Verson & Speed 13
12.Technical Details 14
13.Capacity/Codes 17
14.Interacivity 18
15.Disc Construction 19
16.Hybrid Discs 19
17.Software Standards 22
18.Digital Right Management 24
19.Player Profile 27
20.Variant 27
21.Blu-Ray Recordability 28
22.Bibliography 30

Review Of Litrature

1. Jeremy Reimer (2007-10-07). "New Blu-ray discs with BD+ DRM failing to play on some devices". arstechnica.com. Retrieved on 2007-11-02.
2. Zyber, Joshua (2007-11-23). "High-Def FAQ: Blu-ray Profiles Explained". highdefdigest.com. Retrieved on 2007-12-1
3. Yam, Marcus (2007-01-10). "Three HD Layers Today, Ten Tomorrow" (in English). DailyTech. Retrieved on 2007-04-24.
4. Christian Lysvåg (2008-05-29). "Music on Blu-ray". Music Information Centre Norway. Retrieved on 2008-06-26
5. ^ Joshua Fruhlinger. "First Blu-ray record, Divertimenti, released". engadget. Retrieved on 2008-07-05
6. ^ Zyber, Joshua (2007-11-23). "High-Def FAQ: Blu-ray Profiles Explained". highdefdigest.com. Retrieved on 2007-12-18.

Introduction

Blu-ray Disc (also known as Blu-ray or BD) is an optical disc storage medium. Its main uses are high-definition video and data storage. The disc has the same physical dimensions as standard DVDs and CDs.

The name Blu-ray Disc is derived from the blue laser (violet-colored) used to read and write this type of disc. Because of the beam's shorter wavelength (405 nanometers), substantially more data can be stored on a Blu-ray Disc than on the DVD format, which uses a red (650 nm) laser. A two-layer Blu-ray Disc can store 50 gigabytes, almost six times the capacity of a two-layer DVD, or ten times that of a single-layer DVD.

During the format war over high-definition optical discs, Blu-ray Disc competed with the HD DVD format. On February 19, 2008, Toshiba—the main company supporting HD DVD—announced that it would no longer develop, manufacture, and market HD DVD players and recorders, leading almost all other HD DVD companies to follow suit, effectively ending the format war.

Blu-ray Disc was developed by the Blu-ray Disc Association, a group representing makers of consumer electronics, computer hardware, and motion pictures. As of September 20, 2008, more than 850 Blu-ray Disc titles have been released in the United States and more than 500 Blu-ray Disc titles have been released in Japan. There are expected to be over 1300 Blu-ray Disc titles released in the United States by the end of 2008.

Reverse side of a Blu-ray Disc

Media type : High-density optical disc
Encoding : MPEG-2, H.264/MPEG-4 AVC, and VC-1
Capacity : 25 GB (single layer), 50 GB (dual layer)
Read : 405 nm laser:
Mechanism 1× at 36 Mbit/s
2× at 72 Mbit/s
4× at 144 Mbit/s
6× at 216 Mbit/s
8× at 288 Mbit/s
12× at 432 Mbit/s

Developed by : Blu-ray Disc Association
Usage : Data storage,
High-definition video
High-definition audio
PlayStation 3 games

History

In 1998, commercial HDTV sets began to appear in the consumer market; however, there was no commonly accepted, inexpensive way to record or play HD content. In fact, there was no medium with the storage required to accommodate HD codecs, except JVC's Digital VHS and Sony's HDCAM.Nevertheless, it was well known that using lasers with shorter wavelengths would enable optical storage with higher density. When Shuji Nakamura invented practical blue laser diodes, it was a sensation, although a lengthy patent lawsuit delayed commercial introduction

Origins
Philips and Sony started two projects applying the new diodes: UDO (Ultra Density Optical) and DVR Blue (together with Pioneer), a format of rewritable discs which would eventually become Blu-ray Disc (more specifically, BD-RE).The core technologies of the formats are essentially similar.
The first DVR Blue prototypes were unveiled at the CEATEC exhibition in October 2000.Because the Blu-ray Disc standard places the data recording layer close to the surface of the disc, early discs were susceptible to contamination and scratches and had to be enclosed in plastic cartridges for protection. In February 2002, the project was officially announced as Blu-ray, and the Blu-ray Disc Association was founded by the nine initial members.
The first consumer devices were in stores on April 10, 2003. This device was the Sony BDZ-S77; a BD-RE recorder that was made available only in Japan. The recommended price was US$3800; however, there was no standard for pre-recorded video and no movies were released for this player. The Blu-ray Disc standard was still years away as a newer, more secure DRM system was needed before Hollywood studios would accept it, not wanting to repeat the failure of the Content Scramble System used on DVDs.
Blu-ray Disc Association
The Blu-ray Disc Association (BDA) is the industry consortium that develops and licenses Blu-ray Disc technology and is responsible for establishing format standards and promoting business opportunities for Blu-ray Disc. The BDA is divided into three levels of membership: the Board of Directors, the Contributors, and the General Members.
The "Blu-ray Disc Founder group" was started in May 2002 by nine leading electronic companies: Sony, Matsushita, Pioneer, Philips, Thomson, LG Electronics, Hitachi, Sharp, and Samsung. Spearheaded by Sony Corporation, on February 19, 2002 the companies announced that they were the "Founders" of the Blu-ray Disc and later changed their name to the "Blu-ray Disc Association" on May 18, 2004 to allow more companies to join their development. Some examples of companies that signed in include Apple, TDK, Dell, Hewlett Packard, The Walt Disney Company, Warner Bros. and Universal Music Group. As of December 2007, there are more than 250 members and supporters of the Association.

Members
Board of Directors
The Blu-ray Disc Association website describes the role of the Board of Directors as follows:
"Companies participating in the Board of Directors are active participants of the format creation and key BDA activities. They are selected from the Contributors by election. The board sets an overall strategy and approves key issues. A board member can participate in all activities and attend all meetings. The Blu-ray Disc Founder companies will make up the initial Board of Directors. Annual fee: $ 50,000"
The current 18 board members (as of January 2008) are:[4]
• Apple Inc.
• Dell Inc.
• Hewlett-Packard Company
• Hitachi, Ltd.
• LG Electronics
• Mitsubishi Electric
• Panasonic (Matsushita Electric)
• Pioneer Corporation
• Royal Philips Electronics
• Samsung Electronics
• Sharp Corporation
• Sony Corporation
• Sun Microsystems
• TDK Corporation
• Thomson SA
• Twentieth Century Fox
• Walt Disney Motion Pictures Group / Walt Disney Studios Home Entertainment
• Warner Bros. Entertainment, Inc.

Contributors
Contributors are active participants of the format creation and other key BDA activities. They can be elected to become a member of the Board of Directors. A contributor can attend general meetings and seminars, and can participate in Technical Expert Groups (TEGs), regional Promotion Team activities, and most of the Compliance Committee (CC) activities. Membership requires execution of Contribution Agreement and must be approved by the Board of Directors. Annual fee: $ 20,000

General Membership
"General membership provides access to specific information from Committee discussions. A general member can attend general meetings and seminars. They can participate in specific Regional Promotion Team activities and specific CC activities. Annual fee: $3,000"
Blu-ray Disc format finalized
The Blu-ray Disc physical specifications were finished in 2004. In January 2005, TDK announced that they had developed a hard coating polymer for Blu-ray Discs. The cartridges, no longer necessary, were scrapped. The BD-ROM specifications were finalized in early 2006. AACS LA, a consortium founded in 2004, had been developing the DRM platform that could be used to securely distribute movies to consumers. However, the final AACS standard was delayed, and then delayed again when an important member of the Blu-ray Disc group voiced concerns. At the request of the initial hardware manufacturers, including Toshiba, Pioneer and Samsung, an interim standard was published which did not include some features, like managed copy.
Launch and sales developments
The first BD-ROM players were shipped in the middle of June 2006, though HD DVD players beat them in the race to the market by a few months.
The first Blu-ray Disc titles were released on June 20, 2006. The earliest releases used MPEG-2 video compression, the same method used on DVDs. The first releases using the newer VC-1 and AVC codecs were introduced in September 2006. The first movies using dual layer discs (50 GB) were introduced in October 2006.The first audio-only release was made in March 2008.
The first mass-market Blu-ray Disc rewritable drive for the PC was the BWU-100A, released by Sony on July 18, 2006. It recorded both single and dual layer BD-R as well as BD-RE discs and had a suggested retail price of US$699.
Competition from HD DVD
The DVD Forum (which was chaired by Toshiba) was deeply split over whether to develop the more expensive blue laser technology or not. In March 2002, the forum voted to approve a proposal endorsed by Warner Bros. and other motion picture studios that involved compressing HD content onto dual-layer DVD-9 discs. In spite of this decision, however, the DVD Forum's Steering Committee announced in April that it was pursuing its own blue-laser high-definition solution. In August, Toshiba and NEC announced their competing standard Advanced Optical Disc. It was finally adopted by the DVD Forum and renamed HD DVD the next year, after being voted down twice by Blu-ray Disc Association members, prompting the U.S. Department of Justice to make preliminary investigations into the situation.
HD DVD had a head start in the high definition video market and Blu-ray Disc sales were slow at first. The first Blu-ray Disc player was perceived as expensive and buggy, and there were few titles available. This changed when PlayStation 3 launched, since every PS3 unit also functioned as a Blu-ray Disc player. At CES 2007 Warner proposed Total Hi Def which was a hybrid disc containing Blu-ray on one side and HD DVD on the other but it was never released. By January 2007, Blu-ray discs had outsold HD DVDs, and during the first three quarters of 2007, BD outsold HD DVDs by about two to one. In February 2008, Toshiba withdrew its support for the HD DVD format, leaving Blu-ray as the victor.
Some analysts believe that Sony's PlayStation 3 video game console played an important role in the format war, believing it acted as a catalyst for Blu-ray Disc, as the PlayStation 3 used a Blu-ray Disc drive as its primary information storage medium. They also credited Sony's more thorough and influential marketing campaign. More recently Twentieth Century Fox has cited Blu-ray Disc's adoption of the BD+ anti-copying system as the reason they supported Blu-ray Disc over HD DVD.
Technical specifications
Type Physical size Single layer capacity Dual layer capacity
Standard disc size 12 cm, single sided 25 GB (23.28 GiB) 50 GB (46.56 GiB)
Mini disc size 8 cm, single sided 7.8 GB (7.26 GiB) 15.6 GB (14.53 GiB)

Laser and optics
Blu-ray Disc uses a "blue" (technically violet) laser operating at a wavelength of 405 nm to read and write data. Conventional DVDs and CDs use red and near infrared lasers at 650 nm and 780 nm respectively.
The blue-violet laser's shorter wavelength makes it possible to store more information on a 12 cm CD/DVD sized disc. The minimum "spot size" on which a laser can be focused is limited by diffraction, and depends on the wavelength of the light and the numerical aperture of the lens used to focus it. By decreasing the wavelength, increasing the numerical aperture from 0.60 to 0.85 and making the cover layer thinner to avoid unwanted optical effects, the laser beam can be focused to a smaller spot. This allows more information to be stored in the same area. For Blu-ray Disc, the spot size is 580 nm. In addition to the optical improvements, Blu-ray Discs feature improvements in data encoding that further increase the capacity.
Hard-coating technology
Because the Blu-ray Disc data layer is closer to the surface of the disc, compared to the DVD standard, it was at first more vulnerable to scratches. The first discs were housed in cartridges for protection.
TDK was the first company to develop a working scratch protection coating for Blu-ray Discs. It was named Durabis. In addition, both Sony and Panasonic's replication methods include proprietary hard-coat technologies. Sony's rewritable media are spin-coated with a scratch-resistant and antistatic coating. Verbatim's recordable and rewritable Blu-ray Disc discs use their own proprietary hard-coat technology called Scratch Guard.
Version
BD-RE Versions
There are three versions of rewritable Blu-ray Discs (BD-RE):
Version 1.0
• unique BD File System
• not computer compatible
Version 2.0
• UDF 2.5 file system for computer use
• the use of AACS
• BD-R Version 1.0 follows this specification
Version 3.0
• camcorder (8 cm) discs added
• backward compatible with Version 2.0
• BD-R Version 2.0 follows this specification
BD-R Versions
Version 1.2
• Add BD-R Low to High (BD-R LTH) standard.
Speed
Drive speed Data rate Single layer BD write time
1X 36 Mbit/s 4.5 MB/s 4.39 MiB/s
95 min.
2X 72 Mbit/s 9 MB/s 8.78 MiB/s
47 min.
4X 144 Mbit/s 18 MB/s 17.56 MiB/s
24 min.
6X 216 Mbit/s 27 MB/s 26.34 MiB/s
16 min.
8X 288 Mbit/s 36 MB/s 35.12 MiB/s
12 min.
12X 432 Mbit/s 54 MB/s 52.63 MiB/s
8 min.

BD-R LTH (Low To High)
BD-R LTH is a write-once Blu-ray disc format that features organic dye recording layer. The advantage of BD-R LTH is it can protect a manufacturer's investment in DVD-R/CD-R manufacturing equipment. Only modifications are required to current equipment; no investment in new production lines is required. It is believed this can lower the cost of discs.
Old Blu-ray players and recorders can't utilize BD-R LTH, however, a firmware upgrade could enable the device to access BD-R LTH. Panasonic has already released such a firmware update last November for its DMR-BW200, DMR-BR100 and the MR-BW900/BW800/BW700 models. Pioneer is also expected to ship the first LTH BD drives in Spring of 2008. Furthermore, Sony's PlayStation 3 received firmware upgrade to enable BD-R LTH reading in March, 2008.
Technical details
A Table Comparing the High-definition Optical Media Formats
DVD included for comparison
Mandatory codecs must be supported by the player. Each disc must use one or more of the mandatory codecs.
Blu-ray Disc
HD DVD
DVD

Laser wavelength
405 nm (blue-violet laser) 650 nm (red laser)
Numerical aperture
0.85 0.65 0.6
Storage capacity
(single side) per layer/maximum 25/50 GB 15/30 GB 4.7/8.5 GB
Maximum
bitrate
Raw data transfer 53.95 Mbit/s 36.55 Mbit/s 11.08 Mbit/s
Audio+Video+Subtitles 48.0 Mbit/s 30.24 Mbit/s 10.08 Mbit/s
Video 40.0 Mbit/s 29.4 Mbit/s 9.8 Mbit/s
Mandatory video codecs
H.264/MPEG-4 AVC / VC-1 / MPEG-2
MPEG-1 / MPEG-2

Audio
codecs
lossy
Dolby Digital
Mandatory @ 640 kbit/s Mandatory @ 504 kbit/s Mandatory @ 448 kbit/s
DTS
Mandatory @ 1.5 Mbit/s Optional @ 1.5 Mbit/s
Dolby Digital Plus
Optional @ 1.7 Mbit/s Mandatory @ 3.0 Mbit/s N/A
DTS-HD High Resolution
Optional @ 6.0 Mbit/s Optional @ 3.0 Mbit/s N/A
lossless
Linear PCM
Mandatory Optional
Dolby TrueHD
Optional Mandatory N/A
DTS-HD Master Audio
Optional N/A
Secondary video decoder (PiP)
Mandatory for Bonus View players Mandatory N/A
Secondary audio decoder Mandatory for Bonus View players Mandatory Optional
Interactivity BDMV and Blu-ray Disc Java
Standard Content and Advanced Content
Rudimentary
Internet support Mandatory for BD-Live players Mandatory N/A
Video resolution (maximum)
1920×1080 720×480 (NTSC), 720×576 (PAL)

Frame rates at maximum resolution 24p, 50/60i
25/30p, 50/60i
50/60i

Digital Rights Management
AACS-128bit / BD+ / ROM-Mark
AACS-128bit CSS 40-bit

Region codes
Three area codes
Region free 8 Regions (6 commercial)
Hard coating of disc Mandatory Optional
Capacity/codecs
Blu-ray has a higher maximum disc capacity than HD DVD (50 GB vs. 30 GB for a single sided disc). In September 2007 the DVD Forum approved a preliminary specification for the triple-layer 51GB HD DVD (ROM only) disc though Toshiba never stated whether it was compatible with existing HD DVD players. In September 2006 TDK announced a prototype Blu-ray Disc with a capacity of 200GB. TDK was also the first to develop a Blu-ray prototype with a capacity of 100GB in May 2005. In October 2007 Hitachi developed a Blu-ray prototype with a capacity of 100GB. Hitachi has stated that current Blu-ray drives would only require a few firmware updates in order to play the disc.
The first 50 GB dual-layer Blu-ray Disc release was the movie Click, which was released on October 10, 2006. As of July 2008, over 54% of Blu-ray movies are published on 50 GB dual layer discs with the remainder on 25 GB discs. 85% of HD DVD movies are published on 30 GB dual layer discs, with the remainder on 15 GB discs.
The choice of video compression technology (codec) complicates any comparison of the formats. Blu-ray Disc and HD DVD both support the same three video compression standards: MPEG-2, VC-1 and AVC, each of which exhibits different bitrate/noise-ratio curves, visual impairments/artifacts, and encoder maturity. Initial Blu-ray Disc titles often used MPEG-2 video, which requires the highest average bitrate and thus the most space, to match the picture quality of the other two video codecs. As of July 2008 over 70% of Blu-ray Disc titles have been authored with the newer compression standards: AVC and VC-1. HD DVD titles have used VC-1 and AVC almost exclusively since the format's introduction. Warner Bros., which used to release movies in both formats prior to June 1, 2007, often used the same encode (with VC-1 codec) for both Blu-ray Disc and HD DVD, with identical results. In contrast, Paramount used different encodings: initially MPEG-2 for early Blu-ray Disc releases, VC-1 for early HD DVD releases, and eventually AVC for both formats.
Whilst the two formats support similar audio codecs, their usage varies. Most titles released on the Blu-ray format include Dolby Digital tracks for each language in the region, a DTS-HD Master Audio track for all 20th Century Fox and many upcoming Universal titles, Dolby TrueHD for Disney and Sony Pictures and some Paramount and Warner titles, and for many Blu-ray titles a Linear PCM track for the primary language. On the other hand, most titles released on the HD DVD format include Dolby Digital Plus tracks for each language in the region, and some also include a Dolby TrueHD track for the primary language.
Interactivity
Both Blu-ray Disc and HD DVD have two main options for interactivity (on-screen menus, bonus features, etc.).
Blu-ray's basic mode is known as HDMV or BDMV ("High Definition Movie Mode" or "Blu-ray Disc Movie Mode"), whilst HD DVD's is known as "Standard Content". Both offer modest upgrades from standard DVD, such as the use of more buttons on-screen, a larger colour palette, and expanded programming environment. BDMV is more powerful than Standard Content, and has been used on many Blu-ray disc titles. Standard Content has been used less on HD DVDs. HD DVD's Standard Content is a minor change from standard DVD's sub picture technology, while Blu-ray's BDMV is completely new. This makes transitioning from standard DVD to Standard Content HD DVD relatively simple-- for example, Apple's DVD Studio Pro has supported authoring Standard Content since version 4.0.3. For more advanced interactivity Blu-ray disc supports BD-J while HD DVD supports Advanced Content.

Disc construction
Blu-ray Discs contain their data relatively close to the surface (less than 0.1 mm) which combined with the smaller spot size presents a problem when the surface is scratched as data would be destroyed. To overcome this, TDK, Sony, and Panasonic each have developed a proprietary scratch resistant surface coating. TDK trademarked theirs as Durabis, which has withstood direct abrasion by steel wool and marring with markers in tests.
HD DVD uses traditional material and has the same scratch and surface characteristics of a regular DVD. The data is at the same depth (0.6 mm) as DVD as to minimize damage from scratching. As with DVD the construction of the HD DVD disc allows for a second side of either HD DVD or DVD.
A study performed by Home Media Magazine (August 5, 2007) concluded that HD DVD discs and Blu-ray discs are essentially equal in production cost. Quotes from several disc manufacturers for 25,000 units of HD DVDs and Blu-rays revealed a price differential of only 5-10 cents. (Lowest price: 90 cents versus 100 cents. Highest price: $1.45 versus $1.50.) Another study performed by Wesley Tech (February 9, 2007) arrived at a similar conclusion. Quotes for 10,000 discs show that a 15 gigabyte HD DVD costs $11,500 total, and 25 gigabyte Blu-ray or a 30 gigabyte HD DVD costs $13,000 total. For larger quantities of 100,000 units, the 30 gigabyte HD DVD was more expensive than the 25 gigabyte Blu-ray ($1.55 versus $1.49).
Hybrid discs
At the January 8, 2007 Consumer Electronics Show, Warner Bros. introduced a hybrid technology, Total HD, that would reportedly support both formats on a single disc. The new discs would overlay the Blu-ray and HD DVD layers, placing them respectively 0.1 mm and 0.5 mm beneath the surface. The Blu-ray top layer would act as a two-way mirror, reflecting just enough light for a Blu-ray reader to read and an HD DVD player to ignore. But the following September, Warner President Ron Sanders said that the technology was on hold due to Warner being the only one that would publish on it.
As of January 4, 2008, Warner Bros. announced that they will be supporting the Blu-ray format exclusively after June 1, 2008. This news along with the end of the format war would indicate that the hybrid discs once announced by Warner Bros. will not be put into production.
Copy protection
The primary copy protection system used on both formats is the Advanced Access Content System (AACS). Use of AACS is optional for HD DVD, but mandatory for Blu-ray, which can add thousands of dollars to production costs. Other copy protection systems include:
Blu-ray Disc
HD DVD

• HDCP encrypted digital output
• ROM-Mark watermarking technology (physical layer)
• BD dynamic crypto (BD+)
• HDCP encrypted digital output

Region coding
The Blu-ray specification and all currently available players support region coding. As of July 2008 about 66.7% of Blu-ray Disc titles are region-free and 33.3% use region codes.
The HD DVD specification has no region coding, so an HD DVD disc from anywhere in the world will work in any player. The DVD Forum's steering committee has discussed a request from Disney to add it, but many of the 20 companies on the committee actively oppose it.
Many film titles that are exclusive to Blu-ray in the United States such as Sony's xXx, Fox's Fantastic Four: Rise of the Silver Surfer and Disney's The Prestige, are available on HD DVD in other countries due to different distribution agreements (in fact, The Prestige was released outside the US by once format-neutral studio Warner Bros. Pictures). Because of this, film titles that are exclusive to Blu-ray in the U.S. can be bought on HD DVDs by U.S. consumers purchasing them online from Europe or Asia. Since there is no region coding in HD DVDs, there are no restrictions playing these foreign-bought HD DVDs in an HD DVD player. Similarly, European customers can obtain HD DVD discs from American online retailers for titles that are Blu-ray exclusive or haven't yet been released in their own countries.
Retail price of consumer-writable discs
Disc BD-R
BD-R DL
HD DVD-R
HD DVD-R DL
DVD-R
DVD-R DL

Capacity 25GB 50GB 15GB 30GB 4.7GB 8.5GB
Bulk-Bought Cost $7.99 $21.27 $10.74 $19.85 $0.47 $2.00
Cost Per GB (full disc) $0.44 $0.47 $0.67 $0.62 $0.10 $0.23
Price/Cost Per GB (as of 21.02.2008) $11.99 / $0.48 $36.81 / $0.74 $7 / $0.47 $17.60 / $0.59 $0.47 / $0.10 $2.00 / $0.23
Disc being compared 1 non-rewritable single-layer disc (Verbatim 25GB 2X BD-R) 1 non-rewritable dual-layer disc (Panasonic 50GB 2x BD-R) 1 non-rewritable single-layer disc (Verbatim 15GB 1X HD DVD-R) 1 non-rewritable dual-layer disc (Verbatim 30GB 1X HD DVD-R DL) generic pack of discs generic pack of discs

Software standards
Codecs
Codecs are compression schemes that store audio and video more efficiently, optimizing for either low space usage or quality per megabyte. There are both lossy and lossless compression techniques.
The BD-ROM specification mandates certain codec compatibilities for both hardware decoders (players) and the movie-software (content). For video, all players are required to support MPEG-2, H.264/AVC, and SMPTE VC-1. MPEG-2 is the codec used on regular DVDs, which allows backwards compatibility. H.264/AVC was developed by MPEG and VCEG as a modern successor of H.263 . VC-1 is another MPEG-4 derivative codec mostly developed by Microsoft. BD-ROM titles with video must store video using one of the three mandatory codecs. Multiple codecs on a single title are allowed.
The choice of codecs affects the producer's licensing/royalty costs, as well as the title's maximum runtime, due to differences in compression efficiency. Discs encoded in MPEG-2 video typically limit content producers to around two hours of high-definition content on a single-layer (25 GB) BD-ROM. The more advanced video codecs (VC-1 and H.264) typically achieve a video runtime twice that of MPEG-2, with comparable quality.
MPEG-2 was used by many studios, including Paramount Pictures (which initially used the VC-1 codec for HD DVD releases) for the first series of Blu-ray discs that were launched throughout 2006. Modern releases are now often encoded in either H.264/AVC or VC-1, allowing film studios to place all content on one disc, reducing costs and improving ease of use. Using these codecs will also free many GBs of space for storage of bonus content in HD (1080i/p) as opposed to the SD (480i/p) typically used for most titles. Some studios (such as Warner Bros.) have released bonus content on discs encoded in a different codec than the main feature title; for example the Blu-ray release of Superman Returns uses VC-1 for the feature film and MPEG-2 for bonus content (presumably because it is simply ported from the DVD release).
For audio, BD-ROM players are required to support Dolby Digital, DTS, and linear PCM. Players may optionally support Dolby Digital Plus and DTS-HD High Resolution Audio, as well as lossless formats Dolby TrueHD and DTS-HD Master Audio. BD-ROM titles must use one of the mandatory schemes for the primary soundtrack. A secondary audiotrack, if present, may use any of the mandatory or optional codecs.
For users recording digital television programming, the recordable Blu-ray Disc standard's initial data rate of 36 Mbit/s is more than adequate to record high-definition broadcasts from any source (IPTV, cable/satellite, or terrestrial). BD-Video movies have a maximum data transfer rate of 54 Mbit/s, a maximum AV bitrate of 48 Mbit/s (for both audio and video data), and a maximum video bitrate of 40 Mbit/s. This compares to HD DVD movies which have a maximum data transfer rate of 36 Mbit/s, a maximum AV bitrate of 30.24 Mbit/s, and a maximum video bitrate of 29.4 Mbit/s.

Digital rights management
The Blu-ray Disc format employs several layers of digital rights management.
AACS decryption process
Advanced Access Content System (AACS) is a standard for content distribution and digital rights management. It is developed by AS Licensing Administrator, LLC (AACS LA), a consortium that includes Disney, Intel, Microsoft, Matsushita (Panasonic), Warner Bros., IBM, Toshiba and Sony.
Since appearing in devices in 2006, several successful attacks have been made on the format. The first known attack relied on the trusted client problem. In addition, decryption keys have been extracted from a weakly protected player (WinDVD). Since keys can be revoked in newer releases, this is only a temporary attack and new keys must continually be discovered in order to decrypt the latest discs. This cat-and-mouse game has gone through several cycles and as of August 2008 all current decryption keys are available on the Internet.
BD+ was developed by Cryptography Research Inc. and is based on their concept of Self-Protecting Digital Content. BD+ is effectively a small virtual machine embedded in authorized players. It allows content providers to include executable programs on Blu-ray Discs. Such programs can:
• Examine the host environment, to see if the player has been tampered with. Every licensed playback device manufacturer must provide the BD+ licensing authority with memory footprints that identify their devices.
• Verify that the player's keys have not been changed.
• Execute native code, possibly to patch an otherwise insecure system.
• Transform the audio and video output. Parts of the content will not be viewable without letting the BD+-program unscramble it.
If a playback device manufacturer finds that its devices have been hacked, it can potentially release BD+-code that detects and circumvents the vulnerability. These programs can then be included in all new content releases.
The specifications of the BD+ virtual machine are available only to licensed device manufacturers. A list of licensed adopters is available from the BD+ website.
The first titles using BD+ were released in October 2007. Players from Samsung and LG had problems playing back those titles until the manufacturers updated their firmware, but this problem was later identified as being related to BD-Java use, not BD+.BD+ protection was fully circumvented with the release 6.4.0.0 of Any DVD HD program.
BD-ROM Mark is a small amount of cryptographic data that is stored separately from normal Blu-ray Disc data. Bit-by-bit copies that do not replicate the BD-ROM Mark are impossible to decode. A specially licensed piece of hardware is required to insert the ROM-mark into the media during replication. Through licensing of the special hardware element, the BDA believes that it can eliminate the possibility of mass producing BD-ROMs without authorization.
Player profiles
The BD-ROM specification defines four Blu-ray Disc player profiles which includes an audio only player profile (BD-Audio) that does not require video decoding or BD-J. All three of the video based player profiles (BD-Video) are required to have a full implementation of BD-J, but with varying levels of hardware support.

Feature BD-Audio BD-Video
Grace Period Bonus View BD-Live
Profile 3.0 Profile 1.0 Profile 1.1 Profile 2.0
Built-in persistent memory No 64 KB 64 KB 64 KB
Local storage capability No Optional 256 MB 1 GB
Secondary video decoder (PiP)
No Optional Mandatory Mandatory
Secondary audio decoder No Optional Mandatory Mandatory
Virtual file system
No Optional Mandatory Mandatory
Internet connection capability No No No Mandatory

Variants
Mini Blu-ray Disc
The Mini Blu-ray Disc (also, Mini-BD and Mini Blu-ray) is a compact 8cm (~3in) diameter variant of the Blu-ray Disc that can store approximately 7.5 GB of data. It is similar in concept to the MiniDVD. Recordable (BD-R) and rewritable (BD-RE) versions of Mini Blu-ray Disc have been developed specifically for compact camcorders and other compact recording devices.
BD9/BD5 Blu-ray Disc
BD9 and BD5 are lower capacity variants of the Blu-ray Disc that contain Blu-ray compatible video and audio streams contained on a conventional DVD (650 nm wavelength / red laser) optical disc. Such discs offer the use of the same advanced compression technologies available to Blu-ray discs (including MPEG-4-AVC/H.264, SMPTE-421M/VC-1 and MPEG-2) while using lower cost legacy media. BD9 uses a standard 8152MB DVD9 dual-layer disc while BD5 uses a standard 4482MB DVD5 single-layer disc.
BD9/BD5 discs can be authored using home computers for private showing using standard DVD±R recorders. AACS digital rights management is optional. The BD9/BD5 format was originally proposed by Warner Home Video, as a cost-effective alternative to regular Blu-ray Discs. It was adopted as part of the BD-ROM basic format, file system, and AV specifications. BD9/BD5 is similar to 3× DVD for HD DVD.
AVCREC
AVCREC is an official lower capacity variant of the Blu-ray Disc used for storing Blu-ray Disc compatible content on conventional DVD discs. It is being promoted for use in camcorders, distribution of short HD broadcast content and other cost-sensitive distribution needs. It is similar to HD REC for HD DVD.
Note that AVCREC is not the same as AVCHD content stored on DVD. The latter is a media independent format and is used presently in tape less camcorders that record onto DVD and Blu-ray disks, as well as onto Secure Digital and Memory Stick memory cards. Playing back AVCHD content on a Blu-ray player may require modification of AVCHD directory structure, but does not require re-encoding of video files themselves.
Blu-ray Disc recordable
Blu-ray Disc recordable refers to two optical disc formats that can be recorded with an optical disc recorder. BD-R discs can be written to once, whereas BD-RE can be erased and re-recorded multiple times. The theoretical maximum speed for Blu-ray Discs is about 12× as the speed of rotation (10,000 rpm) causes too much wobbles for the discs to be read properly, similar to the 20× and 52× respective maximum speeds of DVDs and CDs.
Since September 2007, BD-RE was also available in the smaller 8 cm Mini Blu-ray Disc diameter size.
On September 18, 2007, Pioneer and Mitsubishi co-developed BD-R LTH ("Low to High" in groove recording), which features an organic dye recording layer that can be manufactured by modifying existing CD-R and DVD-R production equipment, significantly reducing manufacturing costs.
In February 2008, Taiyo Yuden, Mitsubishi and Maxell released the first BD-R LTH Discs, and in March 2008, Sony's PlayStation 3 gained official support for BD-R LTH Discs with the 2.20 firmware update.
Unlike the previous releases of 120 mm optical discs (i.e. CDs and DVDs), Blu-ray recorders hit the market almost simultaneously with Blu-ray's debut (at least in Japan).

Bibliography

Web-sites:
1.Wikipedia.org
2.Google.co.in
3.Ask.com

communication skills

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:22:00 -0700

Term Paper Of Communication Skills

RICH DAD POOR DAD BY ROBERT.T.KIYOSAKI

Submitted By:- Submitted To:-
Sneha Santosh
BSC-Fashion technology
Roll Nos: 35

CERTIFICATE

This is to certify that the term paper entitled of communication skills completed by Sneha a student of BSC-Fashion Technology, under the guidance of Santosh madam, for the partial fulfillement of the award.
His work has been found……….

Guided by:

ACKNOWLEDGEMENT

Words are not enough to pay gratitude to them who helped me in producing this project. Still I would like to add few words for the people who were a part of this term paper in numerous ways, people who gave unending support right from the stage the idea was conceived.

In particular I wish to thanks our Teacher, SANTOSH, without whose support this project would have been impossible. She has not only helped in giving guidance but also reviewed this project painstaking attention for the details.

I would like to take this opportunity to thanks all the staff members for their unending support which they have provided in many ways.

Last but not the least I would like to thanks all my classmates for overwhelming support through out the making term paper.

SNEHA

PREFACE

It is largely based on Kiyosaki's upbringing and education in Hawaii, although the degree of fictionalization is disputed. Because of the heavy use of allegory, some readers believe that Kiyosaki created Rich Dad as an author surrogate (a literary device), discussed further in the criticism section below. Many readers believe that the "Rich Dad" in the book is actually the founder of Hawaii's widespread ABC Stores.
The book highlights the different attitudes to money, work and life of these two men, and how they in turn influenced key decisions in Kiyosaki's life.
Among some of the book's topics are:
• the value of financial intelligence
• that corporations spend first, then pay taxes, while individuals must pay taxes first
• that corporations are artificial entities that anyone can use, but the poor usually don't know how
According to Kiyosaki and Lechter, wealth is measured as the number of days the income from your assets will sustain you, and financial independence is achieved when your monthly income from assets exceeds your monthly expenses. Each dad had a different way of teaching his son.

ABOUT THE AUTHOR
ROBERT.T.KIYOSAKI
Personal life
A fourth-generation Japanese American, Kiyosaki was born in India and raised in Hawaii . He is the son of the late educator Ralph H. Kiyosaki (1919-1991). After graduating from Hilo High School, he attended the U.S. Merchant Marine Academy in New York, graduating with the class of 1969 as a deck officer. He later served in the Marine Corps as a helicopter gunship pilot during the Vietnam War, where he was awarded the Air Medal. Kiyosaki left the Marine Corps in 1974 and got a job selling copy machines for the Xerox Corporation. In 1977, Kiyosaki started a company that brought to market the first nylon and Velcro "surfer" wallets. The company was moderately successful at first but eventually went bankrupt. In the early 1980s, Kiyosaki started a business that licensed T-shirts for Heavy metal rock bands.[2] Around 1996–1997 he launched Cashflow Technologies, Inc. which operates and owns the Rich Dad (and Cashflow) brand.He is married to Kim Kiyosaki.
Teachings
A large part of Kiyosaki's teachings focus on generating passive income by means of investment opportunities, such as real estate and businesses, with Other Books:
• If you want to be Rich & Happy don't go to School? (1992)
• The Business School for People Who Like Helping People (2001) - endorses multi-level marketing.
• Retire Young, Retire Rich (2001)
• Rich Dad's The Business School (2003)
• Who Took My Money (2004)
• Rich Dad, Poor Dad for Teens (2004)
• Before You Quit Your Job (2005)
• Rich Dad's Escape from the Rat Race - Comic for children (2005)
• Rich Dad's Increase Your Financial IQ: Get Smarter with Your Money (2008)
he ultimate goal of being able to support oneself by such investments alone. In tandem with this, Kiyosaki defines "assets" as things that generate cash inflow, such as rental properties or businesses—and "liabilities" as things that generate cash outflow, such as houses, cars, and so on. Such definitions are somewhat based on the concept of negative gearing. Kiyosaki also argues that financial leverage is critically important in becoming rich.
Kiyosaki stresses what he calls "financial literacy" as the means to obtaining wealth. He says that life skills are often best learned through experience and that there are important lessons not taught in school. He says that formal education is primarily for those seeking to be employees or self-employed individuals, and that this is an "Industrial Age idea." And according to Kiyosaki, in order to obtain financial freedom, one must be either a business owner or an investor, generating passive income.
Kiyosaki speaks often of what he calls "The Cashflow Quadrant," a conceptual tool that aims to describe how all the money in the world is earned. Depicted in a diagram, this concept entails four groupings, split with two lines (one vertical and one horizontal). In each of the four groups there is a letter representing a way in which an individual may earn income
Other Books:
• If you want to be Rich & Happy don't go to School? (1992)
• The Business School for People Who Like Helping People (2001) - endorses multi-level marketing.
• Retire Young, Retire Rich (2001)
• Rich Dad's The Business School (2003)
• Who Took My Money (2004)
• Rich Dad, Poor Dad for Teens (2004)
• Before You Quit Your Job (2005)
• Rich Dad's Escape from the Rat Race - Comic for children (2005)
• Rich Dad's Increase Your Financial IQ: Get Smarter with Your Money (2008)

SUMMARY OF THE BOOK RICH DAD,POOR DAD
BY ROBERT.T.KIYOSAKI

Lesson 1: The Rich Don’t Work For Money
At age 9, Robert Kiyosaki and his best friend Mike asked Mike’s father (Rich Dad) to teach them how to make money. After 3 weeks of dusting cans in one of Rich Dad’s convenience stores at 10 cents a week, Kiyosaki was ready to quit. Rich Dad pointed out this is exactly what his employees sounded like. Some people quit a job because it doesn’t pay well. Others see it as an opportunity to learn something new.
WORK TO LEARN
Next Rich Dad put the two boys to work, this time for nothing. Doing this forced them to think up a source of income, a business scheme. The opportunity came to them upon noticing discarded comic books in the store. The first business plan was hatched. The boys opened a comic book library and employed Mike’s sister at 1$ a week to mind it. Soon they were earning $9.50 a week without having to physically run the library, while kids read as much comics as they could in two
hours after school for only a few cents.

Lesson 2: Why Teach Financial Literacy?
They don’t teach this at school.
T he growing gap between rich and poor is rooted in the antiquated educational system. The system trains people to be good employees, and not employers. The obsolete school system also fails to provide young people with basic financial skills rich people use to grow their wealth.
Know your options and use this knowledge to build a formidable asset column. In an age of instant millionaires it really isn’t about how much money you make, it’s about how much you keep, and how many generations you can keep it.

Lesson 3: Mind Your Own Business
KEEP YOUR DAY JOB BUT START MINDING YOUR OWN BUSINESS.
Kiyosaki sold photocopiers on commission at Xerox. With his earnings he purchased real estate. In 3 years’ time his real estate income was far greater than his earnings at Xerox. He then left the company to mind his own business full time. He knew that in order to get out of the rat race fast, he needed to work harder, sell more copiers and mind his own business.
Don’t spend all your wages. Build a good portfolio of assets and you can spend later when these assets bring you greater income.

Lesson 4: The History of Taxes and the Power of Corporations
Income tax has been levied on citizens in England since 1874. In the United States it was introduced in 1913. Since then what was initially a plan to tax only the rich eventually “trickled down” to the middle class and the poor. The rich have a secret weapon to shelter themselves from heavy taxation. It’s called the Corporation. It isn’t a building with the company name and
logo in brass signage out front. A corporation is simply a legal document in your attorney’s file cabinet duly registered under a government state agency. Corporations offer great tax advantages and protection from lawsuits. It’s the legal way to protect your wealth, and the rich have been using it for generations. Do your own research and find out what taxlaws will bring you the best advantages.
Lesson 5: The Rich Invent Money
Self-confidence coupled with high financial IQ can certainly earn more for you than merely saving a little bit every month.
Make good use of your time and find the best deals.
An example: In the early 90’s the Phoenix economy was bad. Homes once valued at $100,000 sold for $75,000. Kiyosaki shopped at bankruptcy courts and bought the same houses at only $20,000. He resold these properties for $60,000 making a cool $40,000 profit. After six more transactions of the same manner he made a total $190,000 in profit and it only took 30 hours of work time. Rich Dad explains there are Two Types of Investors:

1. Buyers of Packaged Investments.
This is when you call a retail outlet, real estate company, stockbroker
or financial planner and put your money in ready-made investments.
It’s a simple, clean way of investing.
2. The Professional Investor
Design your own investment. Assemble a deal and put together
different components of an opportunity. Rich dad encourages this type.
You need to develop three main skills to be this type of investor
Lesson 6: Work to Learn –Don’t Work for Money
The Author’s Odyssey
After college graduation Robert Kiyosaki joined the Marine Corps. He learned to fly for the love of it. He also learned to lead troops, an important part of management training. His next move was to join Xerox where he learned to overcome his fear of rejection. The thought of knocking on doors and selling copiers terrified him. Soon he was among the top 5 salespeople at the company. For a couple of years he was No.1. Having achieved his objective – overcoming
his shyness and fear—he quit and began minding his own business. Learn skills like PR, marketing, and advertising. Take a second job if it means learning more.

REVIEW OF THE BOOK RICH DAD POOR DAD

In Rich Dad, Poor Dad, Kiyosaki describes the lessons that his two dads taught him about money and its management. To clarify, he had one biological dad and the other was the father of his friend. One of them was highly educated with multiple advanced degrees, the other had an 8th grade education. One was very wealthy, the other regularly struggled with money. Counter-intuitively, the sides were changed on who was wealthy and who was poor. The dad with the 8th grade education, was a wealthy entrepreneur who owned businesses such as restaurants, a construction company and other business ventures. His educated dad spent the majority of life working with very little to show for it.
The first portion of the book is written as a story from the viewpoint of Kiyosaki as a 9 year old kid who learned financial lessons from his rich dad. He performed a number of jobs for him and learned many aspects of business by observing the management, accounting, sales, legal and other aspects. The style of this section was similar to the way The Wealthy Barber was structured in that it teaches financial lessons through narrative style.
A good point Kiyosaki makes is that a house is not an asset though it may be listed this way traditionally. The costs associated with a house such as utilities, property taxes, insurance, and maintenance pull away cash flow. He instead defines an asset as a resource that produces cash. A house actually could be in this category if fully paid for and used as a rental property. (To clarify Kiyosaki does not necessarily recommend buying real estate only with cash. He endorses obtaining financing and taking on debt) I personally think Dave Ramsey's thoughts on this subject of paying cash for investment real estate are more accurate and help to take into account the risk associated with debt.
Other assets could be mutual funds or stocks that generate cash flow as well as intellectual property such as books or music which produce royalties. A business that one owns but doesn't need to be actively involved in the work would also be considered an asset by his definition.
The point he makes is that many people put money into things which do not help to build their wealth and instead cause negative cash flow in some instances through expenses associated with them.
Kiyosaki also promotes a person being creative and figuring out ways to make money in scenarios which might not on the surface look like an opportunity. An example he gives of this is when he worked in a gas station as a kid for very low wages, they sold comic books which were thrown away if not sold by the time the comic salesman returned with the new comics. He collected all of these comics and started a comic book library which charged 10 cents for two hours worth of reading. This allowed kids in the neighborhood to read more comics for the same price that just one would cost. By looking around and finding ways to make money, he identified this opportunity and created a profitable situation.
This philosophy of the book is good in encouraging the building of assets which will continue to increase cash flow as well as the entrepreneurial spirit. One area I do not agree with is the risk level taken on through debt to enable the purchase of real estate. Overall, the book has some good lessons to be gleaned………..

Amorphous Semiconductor Thin Film

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:19:00 -0700

Introduction 1
Prepation 2-6
Application 6-12
Acnowledgement 12-14
References 14

Introduction:-

1950s, Stanford Ovshinsky created an entirely new realm of materials science, which in turn has given new life to the engineering of semiconductors, solar energy, and electric cars.
Stan Ovshinsky was born in Akron, Ohio in 1922. After graduating from high school, he went straight to work. In 1955, he began working the field of amorphous materials, that is, materials that lack a definite crystalline structure. Ovshinsky was the first engineer to devise a method, called "phase change," for crystalizing these disordered materials, with resulting novel uses: for example, films that gain metallic properties without losing their original optical capabilities. One result was amorphous semiconductors --- which the engineering community had previously considered an utter impossibility.
In 1960, Ovshinsky founded Energy Conversion Devices, Inc. (ECD), in order to continue and expand his work in amorphous semiconductors. Meanwhile, engineers nationwide had eagerly entered an entirely new field: "ovonics" (from Ovshinsky Electronics).
Ovshinsky earned numerous patents in the 1970s and '80s for amorphous semiconductor materials. These materials became essential to optoelectronic copying and fax machines, as well as large, flat-panel liquid crystal displays like those of computer monitors. As early as 1970, Ovshinsky had used his ovonic phase change principle to invent a reversible optical memory disk: that is, a prototype rewritable CD-ROM. Today, thirteen high tech companies around the world are developing rewritable CDs using Ovshinsky's technology.
Ovshinsky went on to use his thin-film amorphous silicon to invent a manufacturing method that might do for solar energy what the assembly line did for automobiles. In 1983, he patented a system that allowed photovoltaic solar panels to be manufactured in continuous rolls 1000 feet in length. Ovshinsky's "Continuous Amorphous Solar Cell Production System" operates much like a newspaper rollpress, speedily imprinting a plasma of amorphous silicon semiconductors in a continuous web onto a thin, anodized metal sheet.
The high energy-conversion efficiency of the thin-film cells and the high throughput of the process make Ovshinsky's photovoltaic cells a revolutionary leap forward for solar energy. They have been installed at various sites around and above the globe, from Mexican mountain villages to the Mir space station. Ovshinsky's "Uni-Solar" roofing tiles, for residential buildings, have won Popular Science's "Best of What's New" Grand Award (1996) and Discover Magazine's Discover Award in the Environment category (1997).

PREPARATION:-
1. Process for the preparation of thin semiconductor material films, wherein the process comprises subjecting a semiconductor material wafer having a planar face and whose planeid="DEL-S-00001" date="20070206" ,id="DEL-S-00001" is substantially parallel to a principal crystallographic plane, to the three following stages:
a first stage of implantation by ion bombardment of the face of said wafer by means of ions creating in the volume of said wafer at a depth close to the average penetration depth of said ions, a layer of gaseous microbubbles defining in the volume of said wafer a lower region constituting a majority of the substrate and an upper region constituting the thin id="INS-S-00001" date="20070206" semiconductor material id="INS-S-00001" film, the ions being id="DEL-S-00002" date="20070206" chosen from amongid="DEL-S-00002" hydrogen gas ions id="DEL-S-00003" date="20070206" or rare gas ionsid="DEL-S-00003" andid="INS-S-00002" date="20070206" , wherein id="INS-S-00002" the temperature of the wafer during implantation id="DEL-S-00004" date="20070206" beingid="DEL-S-00004" id="INS-S-00003" date="20070206" is id="INS-S-00003" kept below the temperature at which the gas produced by the implanted ions can escape from the semiconductor by diffusion,
a second stage of intimately contacting the planar face of said wafer with a stiffener constituted by at least one rigid material layer,
a third stage of thermally treating the assembly of said wafer and said stiffener at a temperature above that at which the ion bombardment takes place and adequate to create by a crystalline rearrangement effect in the waferid="INS-S-00004" date="20070206" , a coalescence of hydrogen microbubbles id="INS-S-00004" and a pressure effect in the id="INS-S-00005" date="20070206" hydrogen id="INS-S-00005" microbubbles, a separation between the thin id="INS-S-00006" date="20070206" semiconductor material id="INS-S-00006" film and the majority of the substrate, the stiffener and the planar face of the wafer being kept in intimate contact during said stage.
2. Process for the preparation of thin id="INS-S-00007" date="20070206" semiconductor material id="INS-S-00007" films according to claim 1, wherein the stage of implanting ions in the semiconductor material takes place through one or more layers of materials having a nature and thickness such that they can be traversed by the ions.
3. Process for the id="DEL-S-00005" date="20070206" productionid="DEL-S-00005" id="INS-S-00008" date="20070206" preparation id="INS-S-00008" of thin id="INS-S-00009" date="20070206" semiconductor material id="INS-S-00009" films according to claim 1, wherein the semiconductor comprises a group IV material.
4. Process for the id="DEL-S-00006" date="20070206" productionid="DEL-S-00006" id="INS-S-00010" date="20070206" preparation id="INS-S-00010" of thin id="INS-S-00011" date="20070206" semiconductor material id="INS-S-00011" films id="DEL-S-00007" date="20070206" according to claim 1id="DEL-S-00007" , wherein the id="INS-S-00012" date="20070206" process comprises subjecting a id="INS-S-00012" semiconductor id="DEL-S-00008" date="20070206" isid="DEL-S-00008" id="INS-S-00013" date="20070206" material wafer of id="INS-S-00013" siliconid="DEL-S-00009" date="20070206" ,id="DEL-S-00009" id="INS-S-00014" date="20070206" having a planar face and whose plane is substantially parallel to a principal crystallographic plane, to the three following stages:
a first stage of implantation by ion bombardment of the face of said wafer by means of ions creating in the volume of said wafer at a depth close to the average penetration depth of said ions, a layer of gaseous microbubbles defining in the volume of said wafer a lower region constituting a majority of the substrate and an upper region constituting the thin semiconductor material film, wherein id="INS-S-00014" the implanted id="DEL-S-00010" date="20070206" ion is aid="DEL-S-00010" id="INS-S-00015" date="20070206" ions are id="INS-S-00015" hydrogen gas id="DEL-S-00011" date="20070206" ion,id="DEL-S-00011" id="INS-S-00016" date="20070206" ions and id="INS-S-00016" the wafer temperature during implantation is id="INS-S-00017" date="20070206" kept below the temperature at which the gas produced by the implanted ions can escape from the semiconductor by diffusion and id="INS-S-00017" between 20° and 450° C.id="INS-S-00018" date="20070206" , id="INS-S-00018" id="DEL-S-00012" date="20070206" andid="DEL-S-00012"
id="INS-S-00019" date="20070206" a second stage of intimately contacting the planar face of said wafer with a stiffener constituted by at least one rigid material layer, and
a third stage of thermally treating the assembly of said wafer and said stiffener at a temperature above that at which the ion bombardment takes place and adequate to create by a crystalline rearrangement effect in the wafer and a pressure effect in the microbubbles, a separation between the thin semiconductor material film and the majority of the substrate, the stiffener and the planar face of the wafer being kept in intimate contact during said stage,
wherein id="INS-S-00019" the temperature of the third heat treatment stage exceeds 500° C.
5. Process for the id="DEL-S-00013" date="20070206" productionid="DEL-S-00013" id="INS-S-00020" date="20070206" preparation id="INS-S-00020" of thin id="INS-S-00021" date="20070206" semiconductor material id="INS-S-00021" films according to claim 2, wherein implantation takes place through an encapsulating thermal silicon oxide layer and the stiffener is a silicon wafer covered by at least one silicon oxide layer.
6. Process for the id="DEL-S-00014" date="20070206" productionid="DEL-S-00014" id="INS-S-00022" date="20070206" preparation id="INS-S-00022" of thin id="INS-S-00023" date="20070206" semiconductor material id="INS-S-00023" films according to claim 1, wherein the second stage of intimately contacting the planar face of said wafer with a stiffener takes place by applying an electrostatic pressure.
7. Process for the id="DEL-S-00015" date="20070206" productionid="DEL-S-00015" id="INS-S-00024" date="20070206" preparation id="INS-S-00024" of thin id="INS-S-00025" date="20070206" semiconductor material id="INS-S-00025" films according to claim 1, wherein the stiffener is deposited by one or more methods from within the group consisting of evaporation, sputtering, and chemical vapor deposition with or without plasma assistance or photon assistance.
8. Process for the id="DEL-S-00016" date="20070206" productionid="DEL-S-00016" id="INS-S-00026" date="20070206" preparation id="INS-S-00026" of thin id="INS-S-00027" date="20070206" semiconductor material id="INS-S-00027" films according to claim 1, wherein the stiffener is bonded to said wafer by means of an adhesive substrate.
9. Process for the id="DEL-S-00017" date="20070206" productionid="DEL-S-00017" id="INS-S-00028" date="20070206" preparation id="INS-S-00028" of thin id="INS-S-00029" date="20070206" semiconductor material id="INS-S-00029" films according to claim 1, wherein the stiffener is made to adhere to the wafer by a treatment favoring interatomic bonds.
id="INS-S-00030" date="20070206" 10. Process for the preparation of thin semiconductor material films according to claim 1 further comprising cleaving the thin semiconductor material film from the substrate. id="INS-S-00030"
id="INS-S-00031" date="20070206" 11. Process for the preparation of thin semiconductor material films according to claim 1, wherein the thin semiconductor material films are formed as a continuous film of semiconductor material. id="INS-S-00031"
id="INS-S-00032" date="20070206" 12. Process for the preparation of thin semiconductor material films according to claim 1, wherein the semiconductor material wafer comprises silicon. id="INS-S-00032"
id="INS-S-00033" date="20070206" 13. Process for the preparation of thin semiconductor material films according to claim 1, wherein the semiconductor material wafer comprises germanium. id="INS-S-00033"
id="INS-S-00034" date="20070206" 14. Process for the preparation of thin semiconductor material films according to claim 1, wherein the semiconductor material wafer comprises a silicon-germanium alloy. id="INS-S-00034"
id="INS-S-00035" date="20070206" 15. Process for the preparation of thin semiconductor material films according to claim 1, wherein the semiconductor material wafer comprises silicon carbide. id="INS-S-00035"
id="INS-S-00036" date="20070206" 16. Process for the preparation of thin semiconductor material films according to claim 1, wherein the stiffener comprises a silicon wafer covered by at least one silicon oxide layer. id="INS-S-00036"
id="INS-S-00037" date="20070206" 17. Process for the preparation of thin semiconductor material films, wherein the process comprises subjecting a semiconductor material wafer having a planar face and whose plane is substantially parallel to a principal crystallographic plane, to the three following stages:
a first stage of implantation by hydrogen ion bombardment of the face of said wafer by means of hydrogen ions creating in the volume of said wafer at a depth close to the average penetration depth of said ions, a layer of gaseous microbubbles defining in the volume of said wafer a lower region constituting a majority of the substrate and an upper region constituting the thin semiconductor material film, wherein the temperature of the wafer during implantation is kept below the temperature at which the gas produced by the implanted ions can escape from the semiconductor by diffusion,
a second stage of intimately contacting the planar face of said wafer with a stiffener constituted by at least one rigid material layer,
a third stage of thermally treating the assembly of said wafer and said stiffener at a temperature above that at which the ion bombardment takes place and adequate to create by a crystalline rearrangement effect in the wafer, a coalescence of hydrogen microbubbles and a pressure effect in the hydrogen microbubbles, a separation between the thin semiconductor material film and the majority of the substrate, the stiffener and the planar face of the wafer being kept in intimate contact during said stage. id="INS-S-00037"
id="INS-S-00038" date="20070206" 18. Process for the preparation of thin semiconductor material films according to claim 17, wherein the stage of implanting ions in the semiconductor material takes place through one or more layers of materials having a nature and thickness such that they can be traversed by the ions. id="INS-S-00038"
id="INS-S-00039" date="20070206" 19. Process for the preparation of thin semiconductor material films according to claim 17, wherein the semiconductor material comprises a group IV semiconductor. id="INS-S-00039"
id="INS-S-00040" date="20070206" 20. Process for the preparation of thin semiconductor material films according to claim 17, wherein the semiconductor material wafer comprises silicon.id="INS-S-00040"
id="INS-S-00041" date="20070206" 21. Process for the preparation of thin semiconductor material films according to claim 17, wherein the semiconductor material wafer comprises germanium. id="INS-S-00041"
id="INS-S-00042" date="20070206" 22. Process for the preparation of thin semiconductor material films according to claim 17, wherein the semiconductor material wafer comprises a silicon-germanium alloy. id="INS-S-00042"
id="INS-S-00043" date="20070206" 23. Process for the preparation of thin semiconductor material films according to claim 17, wherein the semiconductor material wafer comprises silicon carbide. id="INS-S-00043"
id="INS-S-00044" date="20070206" 24. Process for the preparation of thin semiconductor material films according to claim 17, wherein implantation takes place through an encapsulating thermal silicon oxide layer. id="INS-S-00044"
id="INS-S-00045" date="20070206" 25. Process for the preparation of thin semiconductor material films according to claim 17, wherein the stiffener comprises a silicon wafer covered by at least one silicon oxide layer. id="INS-S-00045"
id="INS-S-00046" date="20070206" 26. Process for the preparation of thin semiconductor material films according to claim 17, wherein the second stage of intimately contacting the planar face of said wafer with a stiffener takes place by applying an electrostatic pressure. id="INS-S-00046"
id="INS-S-00047" date="20070206" 27. Process for the preparation of thin semiconductor material films according to claim 17, wherein the stiffener is deposited by one or more methods from within the group consisting of evaporation, sputtering, and chemical vapor deposition with or without plasma assistance or photon assistance. id="INS-S-00047"
id="INS-S-00048" date="20070206" 28. Process for the preparation of thin semiconductor material films according to claim 17, wherein the stiffener is bonded to said wafer by means of an adhesive substance. id="INS-S-00048"
id="INS-S-00049" date="20070206" 29. Process for the preparation of thin semiconductor material films according to claim 17, wherein the stiffener is made to adhere to the wafer by a treatment favoring interatomic bonds. id="INS-S-00049"
id="INS-S-00050" date="20070206" 30. Process for the preparation of thin semiconductor material films according to claim 17, which further comprises cleaving the thin semiconductor material film from the substrate. id="INS-S-00050"
id="INS-S-00051" date="20070206" 31. Process for the preparation of thin films according to claim 17, wherein the thin semiconductor material films are formed as a continuous film of semiconductor material. id="INS-S-00051"
id="INS-S-00052" date="20070206" 32. Process for the preparation of thin semiconductor material films, wherein the process comprises subjecting a semiconductor material wafer having a planar face and whose plane is substantially parallel to a principal crystallographic plane, to the three following stages:
a first stage of implantation by ion bombardment of the face of said wafer by means of ions creating in the volume of said wafer at a depth close to the average penetration depth of said ions, a layer of gaseous microbubbles defining in the volume of said wafer a lower region constituting a majority of the substrate and an upper region constituting the thin semiconductor material film, the ions consisting of hydrogen gas ions and, wherein the temperature of the wafer during implantation is kept below the temperature at which the gas produced by the implanted ions can escape from the semiconductor by diffusion,
a second stage of intimately contacting the planar face of said wafer with a stiffener constituted by at least one rigid material layer,
a third stage of thermally treating the assembly of said wafer and said stiffener at a temperature above that at which the ion bombardment takes place and adequate to create by a crystalline rearrangement effect in the wafer, a coalescence of hydrogen microbubbles and a pressure effect in the hydrogen microbubbles, a separation between the thin semiconductor material film and the majority of the substrate, the stiffener and the planar face of the wafer being kept in intimate contact during said stage. id="INS-S-00052"
id="INS-S-00053" date="20070206" 33. Process for the preparation of thin semiconductor material films according to claim 32, wherein the stage of implanting ions in the semiconductor material takes place through one or more layers of materials having a nature and thickness such that they can be traversed by the ions. id="INS-S-00053"
id="INS-S-00054" date="20070206" 34. Process for the preparation of thin semiconductor material films according to claim 32, wherein the semiconductor material comprises a group IV semiconductor. id="INS-S-00054"
id="INS-S-00055" date="20070206" 35. Process for the preparation of thin semiconductor material films according to claim 32, wherein the semiconductor material wafer comrises silicon. id="INS-S-00055"
id="INS-S-00056" date="20070206" 36. Process for the preparation of thin semiconductor material films according to claim 32, wherein the semiconductor material wafer comrises germanium. id="INS-S-00056"
id="INS-S-00057" date="20070206" 37. Process for the preparation of thin semiconductor material films according to claim 32, wherein the semiconductor material wafer comrises a silicon-germanium alloy. id="INS-S-00057"
id="INS-S-00058" date="20070206" 38. Process for the preparation of thin semiconductor material films according to claim 32, wherein the semiconductor material wafer comrises silicon carbide. id="INS-S-00058"
id="INS-S-00067" date="20070206" 47. Process for the preparation of thin semiconductor material films, wherein the process comprises subjecting a semiconductor material wafer having a planar face and whose plane is substantially parallel to a principal crystallographic plane, to the three following stages:
a first stage of implantation by hydrogen ion bombardment of the face of said wafer so as to create in the volume of said wafer at a depth close to the average penetration depth of said ions, a layer of gaseous hydrogen microbubbles defining in the volume of said wafer a lower region constituting a majority of the substrate and an upper region constituting the thin semiconductor material film, wherein the temperature of the wafer during implantation is kept below the temperature at which the gas produced by the implanted ions can escape from the semiconductor by diffusion;
a second stage of intimately contacting the planar face of said wafer with a stiffener constituted by at least one rigid material layer, and
a third stage of thermally treating the assembly of said wafer and said stiffener at a temperature above that at which the ion bombardment takes place and adequate to create by a crystalline rearrangement effect in the wafer and a pressure effect in the microbubbles, a separation between the thin semiconductor material film and the majority of the substrate, the stiffener and the planar face of the wafer being kept in intimate contact during said stage. id="INS-S-00067"
id="INS-S-00068" date="20070206" 48. Process for the preparation of thin semiconductor material films according to claim 47, wherein the stage of implanting ions in the semiconductor material takes place through one or more layers of materials having a nature and thickness such that they can be traversed by the ions. id="INS-S-00068"
id="INS-S-00069" date="20070206" 49. Process for the preparation of thin semiconductor material films according to claim 47, wherein the semiconductor material comprises a group IV semiconductor. id="INS-S-00069"
id="INS-S-00070" date="20070206" 50. Process for the preparation of thin semiconductor material films according to claim 47, wherein the semiconductor material wafer comprises silicon. id="INS-S-00070"
id="INS-S-00071" date="20070206" 51. Process for the preparation of thin semiconductor material films according to claim 47, wherein the semiconductor material wafer comprises germanium. id="INS-S-00071"
id="INS-S-00072" date="20070206" 52. Process for the preparation of thin semiconductor material films according to claim 47, wherein the semiconductor material wafer comprises a silicon-germanium alloy. id="INS-S-00072"
id="INS-S-00073" date="20070206" 53. Process for the preparation of thin semiconductor material films according to claim 47, wherein the semiconductor material wafer comprises silicon carbide. id="INS-S-00073"
id="INS-S-00074" date="20070206" 54. Process for the preparation of thin semiconductor material films according to claim 47, wherein implantation takes place through an encapsulating thermal silicon oxide layer. id="INS-S-00074"
id="INS-S-00075" date="20070206" 55. Process for the preparation of thin semiconductor material films according to claim 47, wherein the stiffener comprises a silicon wafer covered by at least one silicon oxide layer. id="INS-S-00075"
id="INS-S-00076" date="20070206" 56. Process for the preparation of thin semiconductor material films according to claim 47, wherein the second stage of intimately contacting the planar face of said wafer with a stiffener takes place by applying an electrostatic pressure. id="INS-S-00076"
id="INS-S-00077" date="20070206" 57. Process for the preparation of thin semiconductor material films according to claim 47, wherein the stiffener is deposited by one or more methods from within the group consisting of evaporation, sputtering, and chemical vapor deposition with or without plasma assistance or photon assistance. id="INS-S-00077"
id="INS-S-00078" date="20070206" 58. Process for the preparation of thin semiconductor material films according to claim 47, wherein the stiffener is bonded to said wafer by means of an adhesive substance. id="INS-S-00078"
id="INS-S-00079" date="20070206" 59. Process for the preparation of thin semiconductor material films according to claim 47, wherein the stiffener is made to adhere to the wafer by a treatment favoring interatomic bonds. id="INS-S-00079"
id="INS-S-00080" date="20070206" 60. Process for the preparation of thin semiconductor material films according to claim 47, which further comprises cleaving the thin semiconductor material film from the substrate. id="INS-S-00080"
id="INS-S-00081" date="20070206" 61. Process for the preparation of thin films according to claim 47, wherein the thin semiconductor material film is formed as a continuous film of semiconductor material. id="INS-S-00081"
id="INS-S-00082" date="20070206" 62. Process for the preparation of thin semiconductor material films, wherein the process comprises subjecting a semiconductor material wafer having a planar face and whose plane is substantially parallel to a principal crystallographic plane, to the three following stages:
a first stage of implantation by ion bombardment of the face of said wafer by means of hydrogen ions creating, by electronic braking in the wafer, in the volume of said wafer at a depth close to the average penetration depth of said ions, a layer of gaseous hydrogen microbubbles defining in the volume of said wafer a lower region constituting a majority of the substrate and an upper region constituting the thin semiconductor material film, wherein the temperature of the wafer during implantation is kept below the temperature at which the gas produced by the implanted ions can escape from the semiconductor by diffusion;
a second stage of intimately contacting the planar face of said wafer with a stiffener constituted by at least one rigid material layer,
a third stage of thermally treating the assembly of said wafer and said stiffener at a temperature above that at which the ion bombardment takes place and adequate to create by a crystalline rearrangement effect in the wafer and a coalescence of hydrogen microbubbles and a pressure effect in the hydrogen microbubbles, a separation between the thin semiconductor material film and the majority of the substrate, the stiffener and the planar face of the wafer being kept in intimate contact during said stage. id="INS-S-00082"
This thesis deals with the preparation methodologies and the microstructural characteristics concerning semiconductor thin films (including SnO2 thin films, Au/Ge bilayer films, and Pd-Ge alloy thin films) and the metal oxides (including SnO, SnO2, Mn2O3 and Mn3O4 nanocrystals: nanoparticles, nanowires, nanorods, and nanofractals). Firstly, the preparation methodologies and the microstructural characteristics of tin oxides have been investigated in detail and described in chapter 2. This covers the following: (i) the application of x-ray diffraction, scanning electron microscopy, transmission electron microscopy, and high resolution transmission electron microscopy to study tin oxide thin films deposited on Si (100) substrates at room temperature using pulsed laser deposition (PLD) techniques with a sintered cassiterite and subsequently heat-treated tin oxide thin films; (ii) measurement of surface topographies of SnO2 thin films prepared by PLD for various substrate temperatures by scanning electron microscopy, where the concept of fractal geometry was proven useful in describing structures and processes in experimental systems; (iii) preparation of low-dimensional nanostructures of SnO2 thin films with some interesting features of tetragonal rutile structure by PLD, where the as-prepared SnO2 thin films were found to be in the polycrystalline state; (iv) growth of nanocrystalline SnO2 thin films onto glass substrates by PLD, where the thin films were determined to be a polycrystalline SnO2 and an amorphous SnO phase. The nucleation and growth processes of SnO2 nanocrystallites were analysed in detail in order to examine how the PLD technique and operating conditions affect the evolution of grain size, shape, crystallographic characteristics and morphology; (v) experimental and theoretical exploration of quantum dot formation and dynamic scaling behavior of SnO2 nanocrystals in coalescence regime for growth by PLD; (vi) investigation of the mystery of porous SnO2 thin film formation by pulsed delivery based on sintered SnO2 target at room temperature; (vii) preparation of a pure orthorhombic SnO2 thin film by PLD at much lower pressures and temperatures than those of traditional methods; (viii) demonstration that SnO2 nanowires can be synthesized by a PLD process deposited on Si (100) substrates at room temperature; and (ix) synthesis of SnO2 nanorods in bulk quantity by a calcining process based on annealing precursor powders. Secondly, the extended version of metal/semiconductor thin films for the crystallization of amorphous Ge, and the formation of nanocrystals and compounds developed with improved micro- and nanostructured features are described in chapter 3. This chapter includes: (i) investigation of microstructural changes and fractal Ge nanocrystallites in polycrystalline Au/amorphous Ge bilayer films upon annealing by scanning electron microscopy, transmission electron microscopy and x-ray energy-dispersive spectroscopy; (ii) interdiffusion assessment of nanoparticles in fat fractal patterns, where the nanoparticles of polycrystalline Ge have been grown in a freshly cleaved single crystal NaCl (100) substrate, starting from Au/Ge bilayer films prepared by evaporation method during annealing; (iii) investigation of nanocrystal formation and fractal microstructural assessment in Au/Ge bilayer films upon annealing by high-resolution transmission electron microscopy, where the crystallization process was suggested to be a diffusion controlled and random successive nucleation and growth mechanism; (iv) investigation of solid-state reactions and amorphous Ge crystallization for various ratios of thickness (or composition) in Pd-Ge alloy thin films after annealing by transmission electron microscopy; and (v) analysis of grain nucleation, growth and aggregation in Pd-Ge alloy thin films during annealing by fundamental kinetic processes. Thirdly, a novel selective synthesis route for various morphologies of manganese oxides nanocrystals (including nanoparticles, nanorods and nanofractals) and their unique microstructural characteristics are presented in chapter 4. Intricate fundamental properties of manganese oxides nanocrystals are studied. This includes: (i) investigation of the influence of grain size on the vibrational properties of Mn2O3 nanocrystals by Raman and Infrared spectroscopy; (ii) investigation of isothermal grain growth of Mn2O3 nanocrystals at various temperatures between 200 and 500 °C for different annealing times and analyzing the grain growth data using two different models; and (iii) development of a widely applicable chemical reaction route to prepare single-crystal Mn3O4 nanocrystals including nanoparticles, nanorods and nanofractals. The Mn3O4 nanocrystals with tetragonal structure were prepared synchronously by a chemical liquid homogeneous precipitation method, which has been employed to synthesize these nanostructured materials using reactants of MnCl2•4H2O, H2O2, and NaOH under the environment of a suitable surfactant and alkaline solution. To sum up, it is expected that the fabrication methodologies developed and the knowledge of microstructural evolution gained in semiconductor thin films (including SnO2 thin films, Au/Ge bilayer films, and Pd-Ge alloy thin films) and metal oxides (including SnO, SnO2, Mn2O3 and Mn3O4 nanocrystals: nanoparticles, nanowires, nanorods, and nanofractals) will provide an important fundamental basis underpinning further interdisciplinary (physics, chemistry and materials science) research in this field leading to promising exciting opportunities for future technological applications involving these thin film materials.
Semiconductor thin films and their application to dye-sensitized solar cells

Cathodic electrodeposition of titanium dioxide (TiO2) and zinc oxide (ZnO) thin films has been studied in the aim of developing cost-effective alternative routes to the photoelectrode materials for dye-sensitized solar cells (DSCs). Preparation of porous anatase TiO2 thin film modified by cis-dithiocyanato bis(4,4′-dicarboxylic acid-2,2′-bipyridine)ruthenium(II) (N3) dye has been achieved by a three-step process: cathodic electrodeposition of a Ti hydroxide thin film from an acidic aqueous solution containing TiOSO4, H2O2 and KNO3, heat treatment of the film at 400 °C and chemical adsorption of dyes from solution. The photocurrent action spectrum measured at the N3-modified TiO2 thin film electrode in contact with I−/I3− redox electrolyte solution revealed incident photon to current conversion efficiency (IPCE) of 37% in the visible range. While TiO2 needed heat treatment for crystallization, direct electrodeposition of crystalline ZnO was possible from an aqueous solution of Zn(NO3)2. Addition of N3 to the deposition bath made it possible to synthesize porous ZnO/N3 hybrid thin film in one step. IPCE of 24% has been achieved for this film. A sandwich cell using the electrodeposited ZnO/N3 hybrid thin film photoelectrode measured Isc=0.61 mA/cm2, Voc=0.46 V, F.F.=0.46 and η=0.13% under illumination by an artificial light source (500-W Xe lamp equipped with a <420-nm and an IR cutoff filters, INTENSITY=100 mW cm−2), being the first example of a real working DSC fabricated without any heat treatment

ACKNOWLEDGMENTS
We are grateful to Professor Normand Mousseau for sending
us his models of amorphous silicon. We also acknowledge
the support of the National Science Foundation under
Grant Nos. DMR-0074624, DMR-0310933, and DMR-
0205858. P.O. acknowledges support for his research visit to
Ohio University from the Programa de Movilidad de Investigadores
of Ministerio de Educacio´n y Cultura of Spain.
*Electronic address: drabold@helios.phy.ohiou.edu
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11T. Umeda, S. Yamasaki, J. Isoya, and K. Tanaka, Phys. Rev. B 59,
4849 ~1999!.
12P.A. Fedders, D.A. Drabold, P. Ordejon, G. Fabricius, D.
Sanchez-Portal, E. Artcho, and J.M. Soler, Phys. Rev. B 60,
10 594 ~1999!.
13P. Hohenberg and W. Kohn, Phys. Rev. 136, B864 ~1965!.
14M.S. Hybertsen and S.G. Louie, Phys. Rev. B 34, 5390 ~1986!.
15B.G. Pfrommer and S.G. Louie, Phys. Rev. B 58, 12 680 ~1998!.
16 J.-L. Li, G.-M. Rignanese, E.K. Chang, X. Blase, and S.G. Louie,
Phys. Rev. B 66, 035102 ~2002!.
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References :-

-http://lib.cityu.edu.hk/record=b2217892
-www.google .com
-www.answer.com

HARD & SOFT SUPER CONDUCTORS

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:12:00 -0700

Normal electronic conductors have electrical resistance to the motion of electrons whenever a current flows through the material. A voltage must be applied in order to replace this energy lost as heat. A superconductor, however, has no resistance at all. Many metals, but not all, show electrical resistance at ordinary room temperatures but turn superconductive when refrigerated near to absolute zero.
This behaviour of superconductors is exciting today for a variety of commercial applications and in research because the limits of superconductors are a long way from being reached.
. In 1911 superconductivity was first observed in mercury by Dutch physicist Heike Kamerlingh Onnes of Leiden University (shown above). When he cooled it to the temperature of liquid helium, 4 degrees Kelvin (-452F, -269C), its resistance suddenly disappeared. The Kelvin scale represents an "absolute" scale of temperature. Thus, it was necessary for Onnes to come within 4 degrees of the coldest temperature that is theoretically attainable to witness the phenomenon of superconductivity. Later, in 1913, he won a Nobel Prize in physics for his research in this area.
Resistance is classically due to collisions of free electrons with thermally displaced ions with impurities and defects in the metal. This approach can not explain superconductivity as electrons always suffer some collisions so resistance can never be zero. This is put to good use in light bulbs.
The best normal conductors have weak interactions between the electrons and the lattice which is why they are good conductors, but this prevents them from becoming superconductors.
The only way to describe superconductors is to use quantum mechanics. The model used is the BSC theory (named after the 3 men who derived it, Bardeen, Cooper and Schrieffer), which was first suggested in 1957[4]. It states that lattice vibrations play an imp

Type 1 and Type 2 Superconductors
The first superconductors were of little use in a practical sense, because they could not carry a significant amount of current. These are known as type 1 or “soft” superconductors[10]. They require the coldest temperatures (to slow down molecular vibrations sufficiently to allow unimpeded electron flow in accordance with BCS Theory) to become superconductive and exhibit a very sharp transition to a superconducting state and “perfect” diamagnetism.
Diagram 6: Type 1 superconductivity showing a sharp transition
For a type 1 superconductor the critical current is a consequence of the critical magnetic field, Hc[11]. Hc is low in type 1 superconductors along with their critical current densities (important in wire manufacturing) and therefore they have been of little interest to magnet builders or the electric utilities[12].
Type 2 or “hard” superconductors are comprised of metallic compounds and alloys such as “perovskites” (metal oxide ceramics). They achieve higher Tc than type 1 by a mechanism that is still unclear[13]. They differ from type 1 in that their transition from a normal to a superconducting state is gradual across a region of “mixed state” or vortex behaviour. They admit the magnetic field into their interiors while still remaining superconducting. It has been these type 2 superconductors that contemporary scientific and commercial superconducting magnets are wound.
States of Superconductors
In both types of superconductors the electrons combine in pairs under the critical temperature to form macroscopic material waves.
In type 1, the conventional metals and metalloids, the electrons interact with the lattice vibration, whereby both electrons in the Cooper Pair have S= 0 and L=0 (quantum numbers) and can be described by an s-wave function. The wave-function, therefore has the same characteristics along every axis of the lattice.
In type 2 superconductors, the unconventional ceramic compounds, the electron pair processes for an S=0 state; L= 0, 2h, 4h etc. due to quantum mechanic restrictions. The compound always takes on the lowest possible energy and therefore, in most superconductors the (S=0, L=0) state occurs.
7 Paul Brown, Heidelberg University, 2004
High Temperature Superconductors (HTS)
For over 75 years superconductivity remained a low temperature phenomenon, and it was theoretically shown and widely believed that high temperature Superconductivity was impossible, that the highest transition temperature, or critical temperature (Tc), could not go above 30K (according to BCS theory). This changed in 1986, when J.G.Bednorz and K.A.Müller discovered the barium-doped structures of LaCuO4[14], which broke the 30K limit.
With the 30K barrier broken the race was on to find still higher transition temperatures. The first was via strontium substitution: La2-xSrxCuO4 giving a Tc of 38K[15]. It was also found that under extreme pressure the critical temperature could be increased to 50K[16].
The next step was to simulate pressure via chemical substitution. This was done by adding yttrium to the perovskite structure of BaCuO3. Surprisingly the compound (YBa2Cu3O7) went superconducting at 92K[17].
This was important as it put superconductivity in the range of liquid nitrogen and so hundreds of labs joined the race.
It was found that nearly any of the rare earth metals could be substituted for yttrium without any significant change to the transition temperature[18].
The structure of the compound is that of a sandwich with planes of copper oxide in the centre, where the superconducting current flows. The other elements act only as spacers.
The record Tc today is owned by HgBa2Ca2Cu3O8, which by room pressure has a Tc of 135K and under pressure can reach 164K[19]. One of today’s theories predicts an upper limit of 200K for superconductivity, while others predict no limit.
All of these HTSs were brittle ceramic compounds. This is surprising as ceramics are normally insulators. The theory behind this is still not fully understood[20]. This brittleness causes drawbacks in practical applications, such as drawing out wires. Another drawback is the magnetic properties of these materials.
Most HTSs are produced form metastable materials; this means that the thermodynamically stabile reactants are forced into forming the compound either by high pressure and temperature or by doping. This method of synthesis, however, does not represent in any way an absolute criterion for HTS synthesis.
Type II superconductors are, for the most part, comprised of metallic compounds and alloys. This class of superconductors generally has a much higher critical temperature than those in Type I. They achieve a higher critical temperature than Type 1 superconductors by a mechanism that is still not completely understood. It is believed that it relates to the planar layering within the crystalline structure. The highest critical temperature reached is currently 138 K. Debates still arise as to whether or not an upper limit exists for a critical temperature to be found.

REVIEW OF LITERATURE:
1)TITLE:Hard superconductivity of a soft metal in the
quantum regime
MUSTAFA M. ÖZER1, JAMES R. THOMPSON1,2 AND HANNO H. WEITERING
Submitted on:27 january 2006
Superconductivity is inevitably suppressed in reduced dimensionality. The thin superconducting wires or films can be before they lose their
superconducting properties have important technological ramifications and go to the
heart of understanding coherence and robustness of the superconducting state in
quantum-confined geometries. Here, we exploit quantum confinement of itinerant
electrons in a soft metal to stabilize superconductors with lateral dimensions of the
order of a few millimeters and vertical dimensions of only a few atomic layers10.These extremely thin superconductors show no indication of defect- or fluctuationdriven suppression of superconductivity and sustain supercurrents of up to 10% of
the depairing current density. The extreme hardness of the critical state is attributed
to quantum trapping of vortices. This study paints a conceptually appealing, elegant
picture of a model nanoscale superconductor with calculable critical state properties.
It indicates the intriguing possibility of exploiting robust superconductivity at the
nanoscale.
2)Title: Nonlinear diffusion in hard and soft superconductors
Authors: Gilchrist, John; van der Beek, C. J.

Publication Date: 09/1994

Bibliographic Code: 1994PhyC..231..147G
We discuss the diffusion of magnetic flux in a field-cooled (``hard'') superconducting slab in a creep regime in which E ~ |J|σ J. Bryksin and Dorogovtsev recently discussed flux diffusion in a pinningless (``soft'') superconductor in which E ~ |B|J. This problem is closely related to the flux-creep one with σ=1, and provides additional insight into the possible types of behaviour. We list a series of possible long-term asymptotic solutions of a scaling form, which are either analytically exact or accurately calculated. We check numerically that the relevant long-term solution is approached after various initial conditions. Amongst other conclusions we find S=d(In|M|)/d(Int)-->-1/σ or -1/2σ, after application and removal of an additional field, aJump to main content
3)Limited flux jumps in hard superconductors
R G Mints and A L Rakhmanov
Inst. of High Temperatures, Moscow, USSR
Limited flux jumps in superconductors are investigated under the conditions when the heating of the sample is not too high. The surface temperature rise, electric field and magnetic flux change associated with the instability development are calculated. The theory is compared with experiments, and a satisfactory agreement is found.
Print publication: Issue 12 (14 December 1983)
4)Magnetic instabilities in hard superconductors
R G Mints and Aleksandr L Rakhmanov
The magnetic instabilities in hard and combined Type II superconductors in detail give the criteria for stability of the critical state with respect to magnetic-flux jump.Then the total effect of magnetic and thermal diffusion, as well as that of the structure of a combined superconductor, on the magnetic-field value for a flux jump. The theoretical results will be compared with the existing experimental data.
Print publication: Issue 3 (1977)
Superposition of currents in hard superconductors placed into crossed AC and DC magnetic fields
FISHER L. M. (1) ; KALINOV A. V. (1) ; VOLOSHIN I. F. (1) ; BALTAGA I. V. ; IL'ENKO K. V. ; YAMPOL'SKII V. A. ;
Publishing year:1996
The superposition of currents in YBCO melt-textured samples placed into crossed ac and dc magnetic fields is predicted and observed. This superposition is a direct consequence of the critical state model. The dc magnetic field distribution is shown to become uniform wherever the ac field has penetrated. Owing to this nonlinear process, the area of the dc magnetization loop diminishes and eventually disappears completely with an increase of the ac field magnitude. This means that under the action of the external ac field, the static magnetic properties of hard superconductors change and tend to the well known properties of soft ones.

SUMMARY:
Superconductors conduct electricity with little or no resistance. Organic superconductors contain carbon and are less dense than their ceramic or metallic counterparts; they also offer unusual potential for fine-tuning of electrical properties. Argonne National Laboratory long has carried out the major U.S. effort to synthesize and identify organic superconductors. Nearly 100 new superconductors of this type have been produced, with critical temperatures (at which a superconductor loses all electrical resistance) as high as -260 degrees C, or -434 degrees F. Recently, the first superconductor composed entirely of organic components (with no metal atoms) was synthesized, with a transition temperature in this range. Although this remains far lower than the highest known transition temperature for ceramics, scientists still expect that a high-temperature organic superconductor may be possible, such that liquid nitrogen (at -196 degrees C, or -321 degrees F) could be used as the coolant instead of the more costly liquid helium, thus making practical applications more feasible. The new compound has a two-dimensional, layered structure, which may provide significant insight into the nature of superconductivity.
These advances will help scientists develop a theory of how organic superconductors work and contribute to the design of new materials with higher transition temperatures. The all-organic material is ideal for studies of magnetic and charge transport properties because there is no possibility of contamination from metallic impurities.

APPLICATIONS:
Superconductivity already has important applications, such as medical diagnostic equipment, and many more uses are possible if transition temperatures are high enough. The availability of purely organic superconductors greatly expands the possibilities, especially for applications in which weight is a factor
Superconducting high speed train system comprising a rail including at least one elongated hard superconducting member disposed horizontally along the running direction of the train and having a hollow or gap portion extending in the elongated direction, and a train body including a superconducting magnet for generating a magnetic field perpendicular to the hard superconducting member, thereby floating the body from the rail by the magnetic force acting between the superconducting magnet and the hard superconducting member.

LIMITATIONS:

Limitations on performance of Superconductor oversampling ADCs
For the development and optimization of superconductor oversampling modulators, We highlight the importance of specially engineered and parasitic components of the feedback loop. In particular, LR circuits operating as low-pass filters are capable of providing a noticeable SNR improvement and dramatically reducing the dynamic range requirements for used SFQ comparators. On the other hand, the feedback loop delay and time-jitter in timing circuits are able to spoil the potentially extremely high performance of superconductor oversampling ADCs. We also developed a simple formula describing time-jitter in superconductor

BIBLOGRAPHY:
1): arXiv:cond-mat/0601641v1

2)www.iop.org/EJ/abstract/0022-3727/16/12/026
3)http://www.freepatentsonline.com
4)www.sciencedirect.com/science
5)http://physics.aps.org/articles
6) http:/www.msd.anl.go

Ground wave propagation

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:08:00 -0700

Abstract:- I avail this opportunity to convey the entire knowledge of ground wave propagation through this paper.
This paper gives the information about the ground wave propagation, line of sight, various facts effecting ground wave propagation.

Introduction to the problem:- In communication system, there are various methods or techniques for propagation. We use all the methods for communication to one place to another place. Ground wave propagation is one of the most popular technique for transfer the signals of low and medium frequency range of radio waves from one place to another place.

CORE CHAPTER
Radio waves
The word radio means radiation of electromagnetic waves conveying information from one end and receiving such information at other end. Within this meaning such applications telegraphy, telephony, television and a host of navigational age are classified as radio.
All radio waves are electromagnetic waves which travel with speed of light. An electromagnetic wave is created by a local disturbance in the electric and magnetic fields. From its origin, the wave will propagate outwards in all directions. If the medium in which it is propagating (air for example) is the same everywhere, the wave will spread out uniformly in all directions.

The Electromagnetic Spectrum
Frequency Range Band Designation
30-3000 Hz ELF
3-30 kHz VLF
30-300 kHz LF
300-3000 kHz MF
3-30 MHz HF
30-300 MHz VHF
300-3000 MHz UHF
3-30 GHz SHF
30-300 GHz EHF

Propagation of Waves
The process of communication involves the transmission of information from one Location to another. As we have seen, modulation is used to encode the information onto a carrier wave, and may involve analog or digital methods. It is only the characteristics of the carrier wave which determine how the signal will propagate over any significant Distance. This chapter describes the different ways that electromagnetic waves propagate.
Ground Wave
Radio waves in the VLF band propagate in a ground, or surface wave. The wave is connected at one end to the surface of the earth and to the ionosphere at the other. The ionosphere is the region above the troposphere (where the air is), from about 50 to 250 miles above the earth. It is a collection of ions, which are atoms that have some of their electrons stripped off leaving two or more electrically charged objects. The sun's rays cause the ions to form which slowly recombine. The propagation of radio waves in the presence of ions is drastically different than in air, which is why the ionosphere plays an important role in most modes of propagation. Ground waves travel between two limits, the earth and the ionosphere, which acts like a duct. Under normal conditions, the temperature of air gradually decreases with increase in height above ground. When there is a stable high pressure system, a mass of warm air may over run cold air, causing a temperature inversion. Radio waves trapped below the mass can travel great distance with little loss. The area between the earth and the warm air mass is known as duct. In troposphere under normal conditions the air pressure water vapor pressure and temperature reduces with the increase of height above earth. As a result of this the refractive index also reduces with the increase of height. In standard atmosphere the modified refractive index M increases linearly with the increase of height. Under abnormal meteorological conditions particularly under water, variation of dielectric constant of the troposphere with height departs considerably from the standard condition. Thus under certain special conditions the dielectric constant may not decrease at all with height or may even increase with height resulting in the radio wave following the straight line path or curving away from the earth respectively. This results in reduction of the distance to the horizon. Under another type of abnormal conditions the variation of refractive index with height may be much higher than the normal in the region close to the earth. Such abnormal conditions may result when dry air flows from land out over cooler water. Evaporation of moisture from water into the lower layers of air cools the air. Thus the lower layer is cool and rich in moisture while the upper layer is warm and contains less moisture. Thus as the height increases there results increased moisture lapse rate and a temperature inversion that is a rapid increase in temperature instead of reduction. Thus there results rapid decrease in dielectric constant as the wave travels upwards. Typical result of such an abnormal condition is it cause the path of the rays traveling close to the surface of the earth to be bent as much or even more than the curvature of the earth while the rays at greater height are bent less. When the frequency is sufficiently high the region where the variation of refractive index is usually great actually traps energy and causes it to travel along the earth’s surface much as though in a waveguide. This special refraction the electromagnetic wave is called super refraction and the propagation utilizing this super refraction is called duct propagation. The duct propagation may increase the range of space wave communication two or three times the normal line of sight range. Duct propagation becomes possible only at frequencies above one thousand mega hertz. The duct is the region between the upper minimum of the M curve and either the ground or the point where the vertical projection from the upper minimum curve intersects the m curve. When the curve has a negative slope the curvature of the ray is concave downward on a plane earth diagram and the true curvature of the rays is greater than the curvature of the earth. Hence rays which enter the duct with sufficiently small angles are bent until they become horizontal and then are turned downwards.

Since the duct curves with the earth, the ground wave will follow. Therefore very long range propagation is possible using ground waves. Ground wave use for radio communication signal propagation on the long and medium wave bands for local radio communication.

Ground wave propagation
Ground wave propagation is particularly important on the LF and MF portion of the radio spectrum. Ground wave propagation is used to provide relatively local radio communication coverage especially by radio broadcast stations that required to cover particular locality.
Ground wave signal is made up of number of constituent. If the antennas are in the line of sight then there will be a direct wave as well as wave signal.
As the name suggests direct signal is that travels directly between the two antennas and is not effected by the locality. There will also be a reflected signal as the transmission will be reflected by a number of objects including the earth’s surface and any hills, or large buildings that may be present.
In addition to this there is a surface wave. This tends to follow the curvature of the earth and enables coverage to be achieved beyond the horizon. It is the sum of all these components that is known as ground wave.
Beyond the horizon the direct and reflected waves are blocked by the curvature of earth, and the signal is purely made up from the diffracted surface wave. It is for this reason that surface wave is commonly called the ground wave propagation.
The radio signals are spreads out from the transmit along the surface of earth. Instead of just traveling in a straight line the radio signals tend to flow the curvature of the earth. This is because currents are induced in the surface of earth and this action slows down the wave front in this region, causing the wave front of communications signal to tilt downward towards the earth. With the wave front tilted in this direction it is able to curve around the earth and be received well beyond the horizon.

The wave induces currents in the ground over which it passes and thus losses some energy by absorption. This is made up by energy diffracted downward from the upper portion of the wave front.
There is another way in which the surface wave is attenuated. Because of diffraction the wave front gradually tilts over as shown in above figure. As the wave propagates over the earth it tilts more and more, and the increasing tilt causes greater short circuiting of the electric field component of the wave and hence field strength reduction. Eventually, at some distance from the antenna as partly determined by the type of surface over which the ground wave propagates the wave lies down and dies. Thus in the VLF band insufficient range of transmission is cured by increasing the transmitting
power.
Field strength at a distance
Radiation from an antenna by means of the ground wave gives rise to field strength at a distance. If the distance between two antennas is very long the reduction of field strength due to ground and atmospheric absorption reduces the value of voltage received.
Line of Sight
In the VHF band and up, the propagation tends to straighten out into line-of-sight (LOS)
waves. However the frequency is still low enough for some significant effects.
1. Ionospheric scatter. The signal is reflected by the E-region and scattered in all directions. Some energy makes it back to the earth's surface. This seems to be most effective in the range of 600-1000 miles.

Figure 16
1. Tropospheric scatter. Again, the wave is scattered, but this time, by the air itself. This can be visualized like light scattering from fog. This is a strong function of the weather but can produce good performance at ranges under 400 miles.

Figure 17
1. Tropospheric ducting. The wave travels slower in cold dense air than in warm air. Whenever inversion conditions exist, the wave is naturally bent back to the ground. When the refraction matches the curvature of the earth, long ranges can be achieved. This ducting occurs to some extend always and improves the range over true the line-of-sight by about 10 %.
1. Diffraction. When the wave is block by a large object, like a mountain, is can diffract around the object and give coverage where no line-of-sight exists.
Beyond VHF, all the propagation is line-of-sight. Communications are limited by
the visual horizon. The line-of-sight range can be found from the height of the
transmitting and receiving antennas as the addition of four third as the optical horizon.
Effect of frequency
As the wave front of the ground wave surface travels along the earth’s surface it is attenuated. The degree of attenuation is dependent upon a variety of factors. Frequency of the radio signal is one of the major determining factors as losses rise with increasing frequency. As a result it makes this form of propagation impracticable above the bottom end of the HF portion of the spectrum (3 MHz). Typically a signal at 3.0 MHz will suffer an attenuation that may be in the range of 20 to 60 db more than one at 0.5 MHz dependent upon a verity of factors in the signal path including the distance. In view of this it can be seen why high power HF radio broadcast May only audible for a few miles from the transmitting site via the ground wave.
Effect of Ground
The surface wave also very dependent upon the nature of ground over which the signal travels. Ground conductivity, terrain roughness and the dielectric constant all effect the signal attenuation. In addition to this the ground penetration varies, becoming greater at lower frequencies and this means that it is not just the surface conductivity that is of interest. At the higher frequencies this is not of greater importance, but at low frequencies penetration means that ground strata down to 100 meters may have an effect.
Despite all these variables, it is found that terrain with good conductivity gives the best result. Thus soil type and moisture content are of importance. Salty sea water is best, and rich agriculture or marshy land is also good. Dry sandy terrain and city centers are by far the worst. This means sea paths are optimum, although even these are subject to variations through due to roughness of the sea, resulting on path lose being slightly dependent upon the whether it should also be noted that in view of fact that signal penetration has an effect, the water table may have an effect dependent upon the frequency in use.

Effect of polaristion
The type of antenna has a major effect. Vertical polarization is subject to considerably less attenuation than horizontally polarized signals. In some cases the difference can amount to several tens of decibels. It is for the reason that medium wave broadcast stations use vertical antennas, even if they have to be made physically short by adding inductive loading. Ships making use of the mf marine bands often used inverted L antennas as these are able to radiate a significant proportion of the signal that is vertically polarized.

References
• http://www.radio-electronics.com/info/propagation/ground_wave/ground_wave.php

• http://www.articlestreet.com/music/duct-propagation.html

• http://ham-shack.com/propagation.html

Power point

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:07:00 -0700

Dada Motors
Savitri-II, G.T. Road, Near Chowk Dholewal, Ludhiana
Tel: 91-161-2535001-09

Dada Motors
Savitri-III, Near Chowk Dholewal, Ludhiana Tel: 91-161-535001

Dada Motors
Shakti Nagar, Ludhiana
Tel: 91-161-2634670, 2633488

Dada Motors Ltd.
Savitri-I, G.T. Road, Near Chowk Dholewal, Ludhiana
Tel: 91-161-2535001-6

DETAILS AND PRODUCTS OF THE COMPANY.

DADA MOTERS

Dada Motors, a name that has become synonymous with Auto Trade in Northern India, has contemparary history of more than 40 years now.
"Welcome to The World of Dada Motors"
A world of Automobile solutions
Catering to the personal as well as the commerical needs of all the segments, Dada Motors has two to eighteen wheel Automoblile solutions for everyone. Our vast portfolio includes finance, insurance, exchange options and options for genuine spares and standardised services. Our reach is well reflected in our 5,15,000 satisfied customers for new vehicles, over 20,00,000 satisfactorily serviced vehicles and around 30,000 financed cases. These include not only two wheeler finances but also cars, LCV & HCV finance. All these efforts are successfully supported by an active team of 750 employees of the group.
________________________________________A Journey that has been awarded
1962 Joint Family Business handling different authorised dealerships
1972 Dada motors incorporated
1980 Bajaj auto dealership in Ludhiana
1989-All India Best dealer - service award
1989 Bajaj auto dealership in Jalandhar
1991-1992 Largest Bajaj auto seller
1993- Super star dealer award
1993 Telco Dealership
1994-95- Largest selling dealer award for Bajaj
1995- Service incentive award Bajaj
1996- Highest service volume award
1996- Spare part Crore Club award
1996-97- Best service dealer for Tata
1997- Service man of the year award
1998- All India No.1 sales
1999- Highest selling Tata cummins vehicle dealer of North India
1999- 2nd Leading dealer in North India in passanger car divison
2000- All India highest sales award for Bajaj
2001 General Motors dealership at Jalandhar
2001-02 - Award of execellence "market leadership"
2001-02 - Award of execellence "rural marketing"
2001-02 - Award of execellence for segment of geared scooters jalandhar
2001-02 - Award of execellence for segment of 2-wheelers ludhiana
2001-02 - Award of execellence for segment of 2-wheelers mandigobindgarh
2002 Spare part distributor of commercial vehicle of Tata Motors
2002 General Motors Dealership at Amritsar established
2003 General Motors Dealership at Ludhiana established
2003 Dada Insurance established
2003-04 - Dada motors accredited with ISO 9001:2000 “Quality Management System”
2004 Performance award on Tata sales by Mr Ratan Tata (Chairman Tata Motors)
2004 Hindustan Petroleum Petrol Pump dealership at ludhiana
2005 Best performance award in Tata Service by Mr.Ravikant (M.D. Tata Motors)
2005 Best Distributor award in North India for sale process & system by Mr. Ravikant (M.D. Tata Motors)
.
The chairman of the company Mr. Suraj Dada has been awarded with Rashtriya Samman Patra for one of the highest income tax payer for the year 1994-95 to 1998-99.

In addition to this and above all Mr. Suraj Dada, Managing Director of Dada Motors has been awarded with "Dealership of Distinction" award for outstanding contribution to bajaj auto ltd. since 1982.the life time achievement award was presented by Mr. Rahul Bajaj (Chairman and Managing Director of Bajaj Auto. LTD.)

Services Station
S.No Station

1. Dada Motors
G. T. Road, Khajulra, Near Lilly resort
Tel: 91-181-2632030-31

2. Dada Motors
B.M.C Chowk, Jalandhar
Tel: 91-181-2239111

3. Dada Motors Ltd.
G.T. Road (West), Jalandhar Bye Pass, Ludhiana
Tel: 91-161-2741800, 801

Dada Motors Ludhiana H O

Mr Charanjit Singh

+(91)-9915741414

tl1@dadamotors.com

Savitri-1, Nr Dholewal Chowk, Jalandhar City H O, Jalandhar - 144001

Company also make Aluminum and Aluminum Products - Ludhiana, Punjab

ADDRESS:

Dada Inds (Aluminium and Aluminium Products - Ludhiana)

B S F Chowk G T Road, Jalandhar: 144001 Tel. No.2236624

Motors Ludhiana gets highest selling Safari dealer in
Dada Motors Ludhiana gets highest selling Safari dealer in India award Motors Ludhiana gets highest selling Safari dealer in India award

The company was awarded the highest selling Safari dealer in India for the year 2006-2007 by Ratan Tata at the All-India Dealers Meet in Athens, Greece, recently. Dada Motors also pocketed the coveted 'All- India Service Performance Award' at the awards ceremony.

D&B

Title: DADA MOTORS LIMITED (Headquarters)

Price: $4.00
Report Type: D&B Business Overview
Description: Contains basic contact information, industry classifications and names of executives.
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Every day, you make decisions that impact the profitability of your business. Whether these decisions are about managing risk, finding customers or managing your supply base, there is one constant that runs through all of them - data. Dun & Bradstreet's world-class database and expertise can help you cleanse, renew and enhance your current information and keep it up to date.

RMA ELDER NAMED GM DEALER OF THE YEAR FOR THE 6TCONSECUTIVE YEAR
April 29, 2008 - Irma Elder, CEO of Elder Automotive Group in Troy, MI, has been named General Motors' Dealer of the Year and Jack Smith Leadership Award recipient in recognition of her dedication to sales excellence and customer satisfaction for the sixth consecutive year. The award is an exceptional achievement, presented to only 115 of the 6,700 GM dealers nationwide. The award is named after retired GM chairman and program founder Jack Smith.
The Dealer of the Year Award is the highest honor that General Motors bestows each year. The presentation of the award is a testimony to those dealers who exemplify the drive to achieve what GM seeks in all of its authorized dealers. The award recognizes not only sales performance and customer satisfaction, but ongoing operating standards and industry leadership. "I thank General Motors for this award. More importantly, this honor is especially deserved by the people who work for the Elder Automotive Group. Without their extraordinary talents, dedication to their profession and their support of all the important facets of our automobile industry, we could not have achieved such a significant recognition in so many areas of our business," said Elder. "I truly thank them for all that they do so well, and our wonderful customers. Without them we wouldn't be here."
Elder, a native of Mexico, began her career in the automotive industry as an administrative assistant at a Chevrolet dealership. In 1967, she and her husband, Jim, purchased Troy Ford (now Elder Ford). After the death of her husband in 1983, she became the sole owner. Elder now owns 11 dealerships located in Michigan and Florida.
In addition to GM's Dealer of the Year honor, Elder has received numerous awards and accolades for her achievements. They include the GM Mark of Excellence Award, Saab Retail Recognition Award, Automotive Hall of Fame Service Citation Award, Ford Motor Company's President's Award, Latina Pioneer Summit Entrepreneur of the Year, CATCH Hall of Fame Inductee, Visteon Hispanic Network's Adelante Award, Pride of Jaguar Award, Hispanic Person of the Year and the Wonder Woman Crystal Celebration Award.
Elder is a member of the General Motors Minority Dealers Association, the Automotive Women's Association, Detroit Auto Dealers Association (DADA), Metro Detroit Auto Dealers Association, National Auto Dealers Association (NADA), Saab Marketing Vision Team, Christ Child Society, Michigan Hispanic Chamber of Commerce, Troy Chamber of Commerce, Leading Women Entrepreneurs of the World, Project Hope and Women's Forum of Michigan.
Source: General Motors Minority Dealers Association

Chromatography

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:05:00 -0700

CHROMATOGRAPHY

ABSTRACT

Few methods of chemical analysis are truly specific to a particular analyte. It is often found that the analyte of interest must be separated from the myriad of individual compounds that may be present in a sample. As well as providing the analytical scientist with methods of separation, chromatographic techniques can also provide methods of analysis.
Chromatography involves a sample (or sample extract) being dissolved in a mobile phase (which may be a gas, a liquid or a supercritical fluid). The mobile phase is then forced through an immobile, immiscible stationary phase. The phases are chosen such that components of the sample have differing solubilities in each phase. A component which is quite soluble in the stationary phase will take longer to travel through it than a component which is not very soluble in the stationary phase but very soluble in the mobile phase. As a result of these differences in mobilities, sample components will become separated from each other as they travel through the stationary phase.
Techniques such as H.P.L.C. (High Performance Liquid Chromatography) and G.C. (Gas Chromatography) use columns - narrow tubes packed with stationary phase, through which the mobile phase is forced. The sample is transported through the column by continuous addition of mobile phase. This process is called elution. The average rate at which an analyte moves through the column is determined by the time it spends in the mobile phase.

INTRODUCTION

Pictured is a sophisticated gas chromatography system. This instrument records concentrations of acrylonitrile in the air at various points throughout the chemical laboratory.

Chromatography is the collective term for a family of laboratory techniques for the separation of mixtures. It involves passing a mixture dissolved in a "mobile phase" through a stationary phase, which separates the analyte to be measured from other molecules in the mixture and allows it to be isolated.
Chromatography may be preparative or analytical. Preparative chromatography seeks to separate the components of a mixture for further use (and is thus a form of purification). Analytical chromatography normally operates with smaller amounts of material and seeks to measure the relative proportions of analytes in a mixture.

History

The history of chromatography spans from the mid-19th century to the 21st. Chromatography, literally "color writing", was used—and named— in the first decade of the 20th century, primarily for the separation of plant pigments such as chlorophyll. New forms of chromatography developed in the 1930s and 1940s made the technique useful for a wide range of separation processes and chemical analysis tasks, especially in biochemistry.
Some related techniques were developed in the 19th century (and even before), but the first true chromatography is usually attributed to Russian botanist Mikhail Semyonovich Tsvet, who used columns of calcium carbonate for separating plant pigments in the first decade of the 20th century during his research on chlorophyll.
Chromatography began to take its modern form following the work of Archer John Porter Martin and Richard Laurence Millington Synge in the 1940s and 1950s. They laid out the principles and basic techniques of partition chromatography, and their work spurred the rapid development of several lines of chromatography methods: paper chromatography, gas chromatography, and what would become known as high performance liquid chromatography. Since then, the technology has advanced rapidly. Researchers found that the principles underlying Tsvet's chromatography could be applied in many different ways, giving rise to the different varieties of chromatography described below. Simultaneously, advances continually improved the technical performance of chromatography, allowing the separation of increasingly similar molecules
Chromatography terms
• The analyte is the substance that is to be separated during chromatography.
• Analytical chromatography is used to determine the existence and possibly also the concentration of analyte(s) in a sample.
• A bonded phase is a stationary phase that is covalently bonded to the support particles or to the inside wall of the column tubing.
• A chromatogram is the visual output of the chromatograph. In the case of an optimal separation, different peaks or patterns on the chromatogram correspond to different components of the separated mixture.
On the x-axis is the retention time and plotted on the y-axis a signal (for example obtained by a spectrophotometer, mass spectrometer or a variety of other detectors) corresponding to the response created by the analytes exiting the system. In the case of an optimal system the signal is proportional to the concentration of the specific analyte separated
• A chromatograph is equipment that enables a sophisticated separation e.g. gas chromatographic or liquid chromatographic separation.
• Chromatography is a physical method of separation in which the components to be separated are distributed between two phases, one of which is stationary (stationary phase) while the other (the mobile phase) moves in a definite direction.
• The effluent is the mobile phase leaving the column.
• An immobilized phase is a stationary phase which is immobilized on the support particles, or on the inner wall of the column tubing.
• The mobile phase is the phase which moves in a definite direction. It may be a liquid (LC and CEC), a gas (GC), or a supercritical fluid (supercritical-fluid chromatography, SFC). A better definition: The mobile phase consists of the sample being separated/analyzed and the solvent that moves the sample through the column. In one case of HPLC the solvent consists of a carbonate/bicarbonate solution and the sample is the anions being separated. The mobile phase moves through the chromatography column (the stationary phase) where the sample interacts with the stationary phase and is separated.
• Preparative chromatography is used to purify sufficient quantities of a substance for further use, rather than analysis.
• The retention time is the characteristic time it takes for a particular analyte to pass through the system (from the column inlet to the detector) under set conditions.
• The sample is the matter analysed in chromatography. It may consist of a single component or it may be a mixture of components. When the sample is treated in the course of an analysis, the phase or the phases containing the analytes of interest is/are referred to as the sample whereas everything out of interest separated from the sample before or in the course of the analysis is referred to as waste.

The solute refers to the sample components in partition.

Techniques by chromatographic bed shape

Column Chromatography
• The solvent refers to any substance capable of solubilizing other substance, and especially the liquid mobile phase in LC.
• The stationary phase is the substance which is fixed in place for the chromatography procedure. Examples include the silica layer in Chromatography & Thin layer chromatography

Column chromatography is a separation technique in which the stationary bed is within a tube. The particles of the solid stationary phase or the support coated with a liquid stationary phase may fill the whole inside volume of the tube (packed column) or be concentrated on or along the inside tube wall leaving an open, unrestricted path for the mobile phase in the middle part of the tube (open tubular column). Differences in rates of movement through the medium are calculated to different retention times of the sample.
In 1978, W. C. Still introduced a modified version of column chromatography called flash column chromatography (flash). The technique is very similar to the traditional column chromatography, except for that the solvent is driven through the column by applying positive pressure. This allowed most separations to be performed in less than 20 minutes, with improved separations compared to the old method. Modern flash chromatography systems are sold as pre-packed plastic cartridges, and the solvent is pumped through the cartridge. Systems may also be linked with detectors and fraction collectors providing automation. The introduction of gradient pumps resulted in quicker separations and less solvent usage.
A spreadsheet that assists in the successful development of flash columns has been developed. The spreadsheet estimates the retention volume and band volume of analytes, the fraction numbers expected to contain each analyte, and the resolution between adjacent peaks. This information allows users to select optimal parameters for preparative-scale separations before the flash column itself is attempted.
In expanded bed adsorption, a fluidized bed is used, rather than a solid phase made by a packed bed. This allows omission of initial clearing steps such as centrifugation and filtration, for culture broths or slurries of broken cells.

Planar Chromatography

Thin layer chromatography is used to separate components of chlorophyll
Planar chromatography is a separation technique in which the stationary phase is present as or on a plane. The plane can be a paper, serving as such or impregnated by a substance as the stationary bed (paper chromatography) or a layer of solid particles spread on a support such as a glass plate (thin layer chromatography). Different compounds in the sample mixture travel different distances according to how strongly they interact with the stationary phase as compared to the mobile phase . The specific Retardation factor (Rf) of each chemical can be used to aid in the identification of an unknown substance.
Paper Chromatography.
Paper chromatography is a technique that involves placing a small dot of sample solution onto a strip of chromatography paper. The paper is placed in a jar containing a shallow layer of solvent and sealed. As the solvent rises through the paper, it meets the sample mixture which starts to travel up the paper with the solvent. This paper is made of cellulose, a polar substance, and the compounds within the mixture travel farther if they are non-polar. More polar substances bond with the cellulose paper more quickly, and therefore do not travel as far.
Thin layer chromatography

Thin layer chromatography (TLC) is a widely-employed laboratory technique and is similar to paper chromatography. However, instead of using a stationary phase of paper, it involves a stationary phase of a thin layer of adsorbent like silica gel, alumina, or cellulose on a flat, inert substrate. Compared to paper, it has the advantage of faster runs, better separations, and the choice between different adsorbents. For even better resolution and to allow for quantitation, high-performance TLC can be used.

Techniques by physical state of mobile phase

Gas chromatography
Gas chromatography (GC), also sometimes known as Gas-Liquid chromatography, (GLC), is a separation technique in which the mobile phase is a gas. Gas chromatography is always carried out in a column, which is typically "packed" or "capillary" (see below) .
Gas chromatography (GC) is based on a partition equilibrium of analyte between a solid stationary phase (often a liquid silicone-based material) and a mobile gas (most often Helium). The stationary phase is adhered to the inside of a small-diameter glass tube (a capillary column) or a solid matrix inside a larger metal tube (a packed column). It is widely used in analytical chemistry; though the high temperatures used in GC make it unsuitable for high molecular weight biopolymers or proteins (heat will denature them), frequently encountered in biochemistry, it is well suited for use in the petrochemical, environmental monitoring, and industrial chemical fields. It is also used extensively in chemistry research.

Liquid chromatography
Liquid chromatography (LC) is a separation technique in which the mobile phase is a liquid. Liquid chromatography can be carried out either in a column or a plane. Present day liquid chromatography that generally utilizes very small packing particles and a relatively high pressure is referred to as high performance liquid chromatography (HPLC).
In the HPLC technique, the sample is forced through a column that is packed with irregularly or spherically shaped particles or a porous monolithic layer (stationary phase) by a liquid (mobile phase) at high pressure. HPLC is historically divided into two different sub-classes based on the polarity of the mobile and stationary phases. Technique in which the stationary phase is more polar than the mobile phase (e.g. toluene as the mobile phase, silica as the stationary phase) is called normal phase liquid chromatography (NPLC) and the opposite (e.g. water-methanol mixture as the mobile phase and C18 = octadecylsilyl as the stationary phase) is called reversed phase liquid chromatography (RPLC). Ironically the "normal phase" has fewer applications and RPLC is therefore used considerably more.
Specific techniques which come under this broad heading are listed below. It should also be noted that the following techniques can also be considered fast protein liquid chromatography if no pressure is used to drive the mobile phase through the stationary phase. See also Aqueous Normal Phase Chromatography.

Affinity chromatography
Affinity chromatography[4] is based on selective non-covalent interaction between an analyte and specific molecules. It is very specific, but not very robust. It is often used in biochemistry in the purification of proteins bound to tags. These fusion proteins are labelled with compounds such as His-tags, biotin or antigens, which bind to the stationary phase specifically. After purification, some of these tags are usually removed and the pure protein is obtained.

Supercritical fluid chromatography
is commonly used to purify proteins using FPLC. Supercritical fluid chromatography is a separation technique in which the mobile phase is a fluid above and relatively close to its critical temperature and pressure.Ion exchange chromatography utilizes ion exchange mechanism to separate analytes. It is usually performed in columns but the mechanism can be benefited also in planar mode. Ion exchange chromatography uses a charged stationary phase to separate charged compounds including amino acids, peptides, and proteins. In conventional methods the stationary phase is an ion exchange resin that carries charged functional groups which interact with oppositely charged groups of the compound to be retained. Ion exchange chromatography

TECHNIQUES BY SEPARATION MECHANISM

Ion exchange chromatography
Ion-exchange chromatography (or ion chromatography) is a process that allows the separation of ions and polar molecules based on the charge properties of the molecules. It can be used for almost any kind of charged molecule including large proteins, small nucleotides and amino acids. The solution to be injected is usually called a sample, and the individually separated components are called analytes. It is often used in protein purification, water analysis, and quality control.
PRINCIPLE
Stationary phase surface displays ionic functional groups (R-X) that interact with analyte ions of opposite charge. This type of chromatography is further subdivided into cation exchange chromatography and anion exchange chromatography. The ionic compound consisting of the cationic species M+ and the anionic species B- can be retained by the stationary phase.
Cation exchange chromatography retains positively charged cations because the stationary phase displays a negatively charged functional group:
Anion exchange chromatography retains anions using positively charged functional group:Note that the ion strength of either C+ or A- in the mobile phase can be adjusted to shift the equilibrium position and thus retention time.The ion chromatogram shows a typical chromatogram obtained with an anion exchange column.

Size exclusion chromatography
Size exclusion chromatography (SEC) is also known as gel permeation chromatography (GPC) or gel filtration chromatography and separates molecules according to their size (or more accurately according to their hydrodynamic diameter or hydrodynamic volume). Smaller molecules are able to enter the pores of the media and, therefore, take longer to elute, whereas larger molecules are excluded from the pores and elute faster. It is generally a low resolution chromatography technique and thus it is often reserved for the final, "polishing" step of a purification. It is also useful for determining the tertiary structure and quaternary structure of purified proteins, especially since it can be carried out under native solution conditions.

SPECIAL TECHNIQUES

Reversed-phase chromatography
Reversed-phase chromatography is an elution procedure used in liquid chromatography in which the mobile phase is significantly more polar than the stationary phase. Reversed-phase chromatography (RPC) includes any chromatographic method that uses a non-polar stationary phase. The name "reversed phase" has a historical background. In the 1970s most liquid chromatography was done on non-modified silica or alumina with a hydrophilic surface chemistry and a stronger affinity for polar compounds - hence it was considered "normal". The introduction of alkyl chains bonded covalently to the support surface reversed the elution order. Now polar compounds are eluted first while non-polar compounds are retained - hence "reversed phase". All of the mathematical and experimental considerations used in other chromatographic methods apply (ie separation resolution proportional to the column length). Today, reversed-phase column chromatography accounts for the vast majority of analysis performed in liquid chromatography.

Stationary Phases

* Silica Based Stationary Phases
Any inert non-polar substance that achieves sufficient packing can be used for reversed-phase chromatography. The most popular column is a C18 bonded silica (USP classification L1) with 297 columns commercially available [3] This is followed by C8 bonded silica (L7 - 166 columns), pure silica (L3 - 88 columns), cyano bonded silica (L10 - 73 columns) and phenyl bonded silica (L11 - 72 columns). Note that C18, C8 and phenyl are dedicated reversed phase packings while cyano columns can be used in a reversed phase mode depending on analyte and mobile phase conditions. It should be noted at this point that not all C18 columns have identical retention properties. Surface functionalization of silica can be performed in a monomeric or a polymeric reaction with different short-chain organosilanes used in a second step to cover remaining silanol groups (end-capping). While the overall retention mechanism remains the same subtle differences in the surface chemistries of different stationary phases will lead to changes in selectivity.

* Mobile Phase Considerations
Mixtures of water or aqueous buffers and organic solvents are used to elute analytes from a reversed phase column. The solvents have to be miscible with water and the most common organic solvents used are acetonitrile, methanol or tetrahydrofuran (THF). Other solvents can be used such as ethanol, 2-propanol (iso-propyl alcohol). Elution can be performed isocratic (the water-solvent composition does not change during the separation process) or by using a gradient (the water-solvent composition does change during the separation process). The pH of the mobile phase can have an important role on the retention of an analyte and can change the selectivity of certain analytes. Charged analytes can be separated on a reversed phase column by the use of ion-pairing (also called ion-interaction). This technique is known as reversed phase ion-pairing chromatography.

Two-dimensional chromatography
In some cases, the chemistry within a given column can be insufficient to separate some analytes. It is possible to direct a series of unresolved peaks onto a second column with different physico-chemical (Chemical classification) properties. Since the mechanism of retention on this new solid support is different from the first dimensional separation, it can be possible to separate compounds that are indistinguishable by one-dimensional chromatography.

Simulated Moving-Bed Chromatography
In chromatography, the simulated moving bed (SMB) technique is a a variant of high performance liquid chromatography; it is used to separate particles and/or chemical compounds that would be difficult or impossible to resolve otherwise. This increased separation is brought about by a valve-and-column arrangement that is used to lengthen the stationary phase indefinitely.
In the moving bed technique of preparative chromatography the feed entry and the analyte recovery are simultaneous and continuous, but because of practical difficulties with a continuously moving bed in the simulated moving bed technique instead of moving the bed the sample inlet and the analyte exit positions are moved continuously, giving the impression of a moving bed.
True moving bed chromatography (MBC) is only a theoretical concept. Its simulation, SMBC is achieved by the use of a multiplicity of columns in series and a complex valve arrangement, which provides for sample and solvent feed, and also analyte and waste takeoff at appropriate locations of any column, whereby it allows switching at regular intervals the sample entry in one direction, the solvent entry in the opposite direction, whilst changing the analyte and waste takeoff positions appropriately as well.
Ref 3 explains that the advantage of the SMBC is high speed, because a system could be near continuous, whilst its disadvantage is that it only separates binary mixtures. It does not say, but perhaps it can be assumed that this is equivalent with the separation of a single component from a group of compounds. With regard to efficiency it compares with simple chromatography technique like continuous distillation does with batch distillation.

Advantages
When affinity differences between molecules are very small, it is sometimes not possible to improve resolution via mobile- or stationary-phase changes. In these cases, the multi-pass approach of SMB can separate mixtures of those compounds by allowing their small retention time differences to accumulate.
At industrial scale an SMB chromatographic separator is operated continuously, requiring less resin and less solvent than batch chromatography. The continuous operation facilitates operation control and integration into production plants.

Drawbacks
The drawbacks of the SMB are higher investment cost compared to single column operations, a higher complexity, as well as higher maintenance costs. But these drawbacks are effectively compensated by the better yield and a much lower solvent consumption as well as a much higher productivity compared to simple batch separations.
For purifications, in particular the isolation of an intermediate single component or a fraction out of a multicomponent mixture, the SMB is not suited in general. It can only separate two fractions from each other and it does not implement linear solvent gradients as required for the purification of biomolecules.

Applications
In size exclusion chromatography, where the separation process is driven by entropy, it is not possible to increase the resolution attained by a column via temperature or solvent gradients. Consequently, these separations often require SMB, to create usable retention time differences between the molecules or particles being resolved. SMB is also very useful in the pharmaceutical industry, where resolution of molecules having different chirality must be done on a very large scale.
For the production of Fructose e.g. in High fructose corn syrup or amino-acids, biological-acids, etc. industrial scale chromatography is used.

Fast protein liquid chromatography

Fast protein liquid chromatography (FPLC) is a term applied to several chromatography techniques which are used to purify proteins. Many of these techniques are identical to those carried out under high performance liquid chromatography, however use of FPLC techniques are typically for preparing large scale batches of a purified product.Basically Fast Protein Liquid Chromatography, usually referred to as FPLC, is a form of column chromatography used to separate or purify proteins from complex mixtures. It is very commonly used in biochemistry and enzymology. Columns used with an FPLC can separate macromolecules based on size, charge distribution (ion exchange), hydrophobicity, or biorecognition (as with affinity chromatography).
The system setup is very similar to that of an HPLC, although the materials, buffers, and pressures used are usually different.
Technically, FPLC is the trade name for the protein chromatography system developed by Pharmacia (now GE Healthcare), and is now sold under the ÄKTA brand. However, it is often used as a genericized trademark to describe high-pressure chromatography purification of proteins.

Countercurrent chromatography

Countercurrent chromatography (CCC) or partition chromatography is a category of liquid-liquid chromatography techniques.[1] Chromatography in general is used to separate components of a mixture based on their differing affinities for mobile and stationary phases of a column. The components can then be analyzed separately by various sorts of detectors which may or may not be integrated into an apparatus. Partition chromatography is based on differences in capacity factor,k, and distribution coefficient,Kd.of the analytes using liquid stationary and mobile liquid phase.In liquid-liquid chromatography, both the mobile and stationary phases are liquid. In contrast, standard column chromatography uses a solid stationary phase and a liquid mobile phase, while gas chromatography uses a liquid stationary phase on a solid support and a gaseous mobile phase. By eliminating solid supports, permanent adsorption of the analyte onto the column is avoided, and a near 100% recovery of the analyte can be achieved. The instrument is also easily switched between various modes of operation simply by changing solvents. With liquid-liquid chromatography, researchers are not limited by the composition of the columns commercially available for their instrument. Nearly any pair of immiscible solutions can be used in liquid-liquid chromatography, and most instruments can be operated in standard or reverse-phase modes. Solvent costs are also generally cheaper than for HPLC, and the cost of purchasing and disposing of solid adsorbents is completely eliminated. Another advantage is that experiments conducted in the lab can easily be scaled to industrial volumes interface between them has a large area, and the analyte can move between the phases according to its partition coefficient. A partition coefficient is a ratio of the amount of analyte found in each of the solvents at equilibrium and is related to the analyte's affinity for one over the other. The mobile phase is mixing with then settling from the stationary phase throughout the column. The degree of stationary phase retention (inversely proportional to the amount of stationary phase loss or "bleed" in the course of a separation) is a crucial parameter. Higher quality instruments have greater stationary phase retention. The settling time is a property of the solvent system and the sample matrix, both of which. When GC or HPLC is done with large volumes, resolution is lost due to issues with surface-to-volume ratios and flow dynamics; this is avoided when both phases are liquid.
CCC can be thought of as occurring in three stages: mixing, settling, and separation (although they often occur continuously). Mixing of the phases is necessary so that the greatly influence stationary phase retention

Chiral chromatography
Chiral chromatography involves the separation of stereoisomers. In the case of enantiomers, these have no chemical or physical differences apart from being three dimensional mirror images. Conventional chromatography or other separation processes are incapable of separating them. To enable chiral separations to take place, either the mobile phase or the stationary phase must themselves be made chiral, giving differing affinities between the analytes. Chiral chromatography HPLC columns (with a chiral stationary phase) in both normal and reversed phase are commercially available.

CONCLUSIONS

In any chemical or bioprocessing industry, the need to separate and purify a product from a complex mixture is a necessary and important step in the production line. Today, there exists a wide market of methods in which industries can accomplish these goals. Chromatography is a very special separation process for a multitude of reasons! First of all, it can separate complex mixtures with great precision. Even very similar components, such as proteins that may only vary by a single amino acid, can be separated with chromatography. In fact, chromatography can purify basically any soluble or volatile substance if the right adsorbent material, carrier fluid, and operating conditions are employed. Second, chromatography can be used to separate delicate products since the conditions under which it is performed are not typically severe. For these reasons, chromatography is quite well suited to a variety of uses in the field of biotechnology, such as separating mixtures of proteins

REFERENCES

1. http://en.wikipedia.org/wiki/Chromatography
2. http://www.rpi.edu/dept/chem-eng/Biotech-Environ/CHROMO/chromintro.html
3.http://teaching.shu.ac.uk/hwb/chemistry/tutorials/chrom/chrom1.htm

BOHR’S MODEL OF ATOM

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:03:00 -0700

ACKNOWLEDGEMENT

“THE BEST USE OF PEN IS FOR EXCELLENCE”. With these words I RAJINDER KAUR I would like to thank the almighty GOD for showing all his blessings and giving me strength for completing this project. He is the without whom this all would never had been possible. I would also like to my teacher Prof J.P. Aggarwal for providing me with this opportunity. The Vice Chancellor of LOVELY PROFESSIONAL UNIVERSITY Dr. Vijay Gupta for providing his encouraging words. Lastly my parents without whom my life would never had been possible and also making me part for this esteemed institution.

RAJINDER KAUR

INDEX

S.No TOPIC PGE NO.
1 Acknowledgement 2
2 Bohr s model 3
3 Shortcomings 9
4 Bibliography 11

Bohr model

The Rutherford-Bohr model of the hydrogen atom (Z = 1) or a hydrogen-like ion (Z > 1), where the negatively charged electron confined to an atomic shell encircles a small positively charged atomic nucleus, and an electron jump between orbits is accompanied by an emitted or absorbed amount of electromagnetic energy hν. The orbits that the electron may travel in are shown as grey circles; their radius increases as n2, where n is the principal quantum number. The transition depicted here produces the first line of the Balmer series, and for hydrogen (Z = 1) results in a photon of wavelength 656 nm (red).
In atomic physics, the Bohr model created by Niels Bohr depicts the atom as a small, positively charged nucleus surrounded by electrons that travel in circular orbits around the nucleus—similar in structure to the solar system, but with electrostatic forces providing attraction, rather than gravity. This was an improvement on the earlier cubic model (1902), the plum-pudding model (1904), the Saturnian model (1904), and the Rutherford model (1911). Since the Bohr model is a quantum physics-based modification of the Rutherford model, many sources combine the two, referring to the Rutherford-Bohr model.
Introduced by Niels Bohr in 1913, the model's key success lay in explaining the Rydberg formula for the spectral emission lines of atomic hydrogen; while the Rydberg formula had been known experimentally, it did not gain a theoretical underpinning until the Bohr model was introduced. Not only did the Bohr model explain the reason for the structure of the Rydberg formula, but it provided a justification for its empirical results in terms of fundamental physical constants.
The Bohr model is a primitive model of the hydrogen atom. As a theory, it can be derived as a first-order approximation of the hydrogen atom using the broader and much more accurate quantum mechanics, and thus may be considered to be an obsolete scientific theory. However, because of its simplicity, and its correct results for selected systems (see below for application), the Bohr model is still commonly taught to introduce students to quantum mechanics, before moving on to the more accurate but more complex valence shell atom. A related model was originally proposed by Arthur Erich Haas in 1910, but was rejected. The quantum theory of the period between Planck's discovery of the quantum (1900) and the advent of a full-blown quantum mechanics (1925) is often referred to as the old quantum theory.

Origin
In the early 20th century, experiments by Ernest Rutherford established that atoms consisted of a diffuse cloud of negatively charged electrons surrounding a small, dense, positively charged nucleus. Given this experimental data, Rutherford naturally considered a planetary-model atom, the Rutherford model of 1911 — electrons orbiting a solar nucleus — however, said planetary-model atom has a technical difficulty. The laws of classical mechanics (i.e. the Larmor formula), predict that the electron will release electromagnetic radiation while orbiting a nucleus. Because the electron would lose energy, it would gradually spiral inwards, collapsing into the nucleus. This atom model is disastrous, because it predicts that all matter is unstable.
Also, as the electron spirals inward, the emission would gradually increase in frequency as the orbit got smaller and faster. This would produce a continuous smear, in frequency, of electromagnetic radiation. However, late 19th century experiments with electric discharges through various low-pressure gasses in evacuated glass tubes had shown that atoms will only emit light (that is, electromagnetic radiation) at certain discrete frequencies.
To overcome this difficulty, Niels Bohr proposed, in 1913, what is now called the Bohr model of the atom. He suggested that electrons could only have certain classical motions:
The electrons can only travel in special orbits: at a certain discrete set of distances from the nucleus with specific energies.
The electrons do not continuously lose energy as they travel. They can only gain and lose energy by jumping from one allowed orbit to another, absorbing or emitting electromagnetic radiation with a frequency ν determined by the energy difference ΔE = E2 − E1 of the levels according to Bohr's formula where h is Planck's constant.
the frequency of the radiation emitted at an orbit with period T is as it would be in classical mechanics--- it is the reciprocal of the classical orbit period:
The significance of the Bohr model is that the laws of classical mechanics apply to the motion of the electron about the nucleus only when restricted by a quantum rule. Although rule 3 is not completely well defined for small orbits, because the emission process involves two orbits with two different periods, Bohr could determine the energy spacing between levels using rule 3 and come to an exactly correct quantum rule: the angular momentum L is restricted to be an integer multiple of a fixed unit:
where n = 1,2,3,… and is called the principal quantum number. The lowest value of n is 1. This gives a smallest possible orbital radius of 0.0529 nm. This is known as the Bohr radius. Once an electron is in this lowest orbit, it can get no closer to the proton. Starting from the angular momentum quantum rule Bohr[1] was able to calculate the energies of the allowed orbits of the hydrogen atom and other hydrogenlike atoms and ions.

Electron energy levels
The Bohr model gives almost exact results only for a system where two charged points orbit each other at speeds much less than that of light. This not only includes one-electron systems such as the hydrogen atom, singly-ionized helium, doubly ionized lithium, but it includes positronium and Rydberg states of any atom where one electron is far away from everything else. It can be used for K-line X-ray transition calculations if other assumptions are added (see Moseley's law below). In high energy physics, it can be used to calculate the masses of heavy quark mesons.
To calculate the orbits requires two assumptions:
1. Classical mechanics
The electron is held in a circular orbit by electrostatic attraction. The centripetal force is equal to the Coulomb force.
where me is the mass and e is the charge of the electron. This determines the speed at any radius:
It also determines the total energy at any radius:
The total energy is negative and inversely proportional to r. This means that it takes energy to pull the orbiting electron away from the proton. For infinite values of r, the energy is zero, corresponding to a motionless electron infinitely far from the proton. The total energy is half the potential energy, which is true for non circular orbits too by the virial theorem.
For larger nuclei, the only change is that ke2 is everywhere replaced by Zke2 where Z is the number of protons. For positronium, me is replaced by the reduced mass me / 2.
2. Quantum rule
The angular momentum is an integer multiple of :
, in which .
Substituting the expression for the velocity gives an equation for r in terms of n:
so that the allowed orbit radius at any n is:
The smallest possible value of r

is called the Bohr radius.
The energy of the n-th level is determined by the radius:
So an electron in the lowest energy level of hydrogen (n = 1) has 13.606 eV less energy than a motionless electron infinitely far from the nucleus. The next energy level at (n = 2) is -3.4 eV. The third (n = 3) is -1.51 eV, and so on. For larger values of n, these are also the binding energies of a highly excited atom with one electron in a large circular orbit around the rest of the atom.
The combination of natural constants in the energy formula is called the Rydberg energy RE:
This expression is clarified by interpreting it in combinations which form more natural units:

: the rest energy of the electron (= 511 keV)
: the fine structure constant
For nuclei with Z protons, the energy levels are:
(Heavy Nuclei)
When Z is approximately 137 (about 1/α), the motion becomes highly relativistic. Then the Z2 cancels the α2 in R, so the orbit energy begins to be comparable to rest energy. Sufficiently large nuclei, if they were stable, would reduce their charge by creating a bound electron from the vacuum, ejecting the positron to infinity. This is the theoretical phenomenon of electromagnetic charge screening which predicts a maximum nuclear charge. Emission of such positrons has been observed in the collisions of heavy ions to create temporary super-heavy nuclei[citation needed].
For positronium, the formula uses the reduced mass. For any value of the radius, the electron and the positron are each moving at half the speed around their common center of mass, and each has only one fourth the kinetic energy. The total kinetic energy is half what it would be for a single electron moving around a heavy nucleus.
(Positronium)
Rydberg formula
The Rydberg formula, which was known empirically before Bohr's formula, is now in Bohr's theory seen as describing the energies of transitions or quantum jumps between one orbital energy level, and another. Bohr's formula gives the numerical value of the already-known and measured Rydberg's constant, but now in terms of more fundamental constants of nature, including the electron's charge and Planck's constant.
When the electron moves from one energy level to another, a photon is emitted. Using the derived formula for the different 'energy' levels of hydrogen one may determine the 'wavelengths' of light that a hydrogen atom can emit.
The energy of a photon emitted by a hydrogen atom is given by the difference of two hydrogen energy levels:
where nf is the final energy level, and ni is the initial energy level.
Since the energy of a photon is
This is known as the Rydberg formula, and the Rydberg constant R is RE / hc, or RE / 2π in natural units. This formula was known in the nineteenth century to scientists studying spectroscopy, but there was no theoretical explanation for this form or a theoretical prediction for the value of R, until Bohr. In fact, Bohr's derivation of the Rydberg constant, as well as the concomitant agreement of Bohr's formula with experimentally observed spectral lines of the Lyman (nf = 1), Balmer (nf = 2), and Paschen (nf = 3) series, and successful theoretical prediction of other lines not yet observed, was one reason that his model was immediately accepted.
Shell model of the atom
Bohr extended the model of Hydrogen to give an approximate model for heavier atoms. This gave a physical picture which reproduced many known atomic properties for the first time.
Heavier atoms have more protons in the nucleus, and more electrons to cancel the charge. Bohr's idea was that each discrete orbit could only hold a certain number of electrons. After that orbit is full, the next level would have to be used. This gives the atom a shell structure, in which each shell corresponds to a Bohr orbit.
This model is even more approximate than the model of hydrogen, because it treats the electrons in each shell as non-interacting. But the repulsions of electrons is taken into account somewhat by the phenomenon of screening. The electrons in outer orbits do not only orbit the nucleus, but they also orbit the inner electrons, so the effective charge Z that they feel is reduced by the number of the electrons in the inner orbit.

For example, the lithium atom has two electrons in the lowest 1S orbit, and these orbit at Z=2. Each one sees the nuclear charge of Z=3 minus the screening effect of the other, which crudely reduces the nuclear charge by 1 unit. This means that the innermost electrons orbit at approximately 1/4th the Bohr radius. The outermost electron in lithium orbits at roughly Z=1, since the two inner electrons reduce the nuclear charge by 2. This outer electron should be at nearly one Bohr radius from the nucleus. Because the electrons strongly repel each other, the effective charge description is very approximate, the effective charge Z doesn't usually come out to be an integer. But Moseley's law experimentally probes the innermost pair of electrons, and shows that they do see a nuclear charge of approximately Z-1, while the outermost electron in an atom or ion with only one electron in the outermost shell orbits a core with effective charge Z-k where k is the total number of electrons in the inner shells.
The shell model was able to qualitatively explain many of the mysterious properties of atoms which became codified in the late 19th century in the periodic table of the elements. One property was the size of atoms, which could be determined approximately by measuring the viscosity of gases and density of pure crystalline solids. Atoms tend to get smaller as you move to the right in the periodic table, becoming much bigger at the next line of the table. Atoms to the right of the table tend to gain electrons, while atoms to the left tend to lose them. Every element on the last column of the table is chemically inert (noble gas).
In the shell model, this phenomenon is explained by shell-filling. Successive atoms get smaller because they are filling orbits of the same size, until the orbit is full, at which point the next atom in the table has a loosely bound outer electron, causing it to expand. The first Bohr orbit is filled when it has two electrons, and this explains why helium is inert. The second orbit allows eight electrons, and when it is full the atom is neon, again inert. The third orbital contains eight again, except that in the more correct Sommerfeld treatment (reproduced in modern quantum mechanics) there are extra "d" electrons. The third orbit may hold an extra 10 d electrons, but these positions are not filled until a few more orbitals from the next level are filled (Filling the n=3 d orbitals produces the 10 transition elements). The irregular filling pattern is an effect of interactions between electrons, which are not taken into account in either the Bohr or Sommerfeld models, and which are difficult to calculate even in the modern treatment.

Shortcomings
The Bohr model gives an incorrect value for the ground state orbital angular momentum. The angular momentum in the true ground state is known to be zero. Although mental pictures fail somewhat at these levels of scale, an electron in the lowest modern "orbital" with no orbital momentum, may be thought of as not to rotate "around" the nucleus at all, but merely to go tightly around it in an ellipse with zero area (this may be pictured as "back and forth", without striking or interacting with the nucleus). This is only reproduced in a more sophisticated semiclassical treatment like Sommerfeld's. Still, even the most sophisticated semiclassical model fails to explain the fact that the lowest energy state is spherically symmetric--- it doesn't point in any particular direction.
In modern quantum mechanics, the electron in hydrogen is a spherical cloud of probability which grows denser near the nucleus. The rate-constant of probability-decay in hydrogen is equal to the inverse of the Bohr radius, but since Bohr worked with circular orbits, not zero area ellipses, the fact that these two numbers exactly agree, is considered a "coincidence." (Though many such coincidental agreements are found between the semi-classical vs. full quantum mechanical treatment of the atom; these include identical energy levels in the hydrogen atom, and the derivation of a fine structure constant, which arises from the relativistic Bohr-Sommerfield model (see below), and which happens to be equal to an entirely different concept, in full modern quantum mechanics).
The Bohr model also has difficulty with, or else fails to explain:
Much of the spectra of larger atoms. At best, it can make predictions about the K-alpha and some L-alpha X-ray emission spectra for larger atoms, if two additional ad hoc assumptions are made (see Moseley's law above). Emission spectra for atoms with a single outer-shell electron (atoms in the lithium group) can also be approximately predicted. Also, if the empiric electron-nuclear screening factors for many atoms are known, many other spectral lines can be deduced from the information, in similar atoms of differing elements, via the Ritz-Rydberg combination principles (see Rydberg formula). All these techniques essentially make use of Bohr's Newtonian energy-potential picture of the atom.
The theory does not predict the relative intensities of spectral lines; although in some simple cases, Bohr's formula or modifications of it, was able to provide reasonable estimates (for example, calculations by Kramers for the Stark effect).
The existence of fine structure and hyperfine structure in spectral lines, which are known to be due to a variety of relativistic and subtle effects, as well as complications from electron spin.

BIBLEOGRAPHY

• www.google.com
• www.ieee.com
• www.wikipedia.com
• www.phyforums.com
• APPLIED CHEMISTRY

DETERMINATION OF WAVE FUNCTION OF sp² & sp³ hybrid orbital

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:03:00 -0700

INTRODUCTION:- In Hybridization , there are three types of hybrid orbital – sp, sp², sp³ .Subscripts 2, 3 indicates that the shapes of hybrid orbitals are different from those of the parent orbitals. e.g.:- s orbital is spherical & p orbital is dumb-bell shaped and thus, the shape of the bigger lobe of sp hybrid orbital which results from the mixing of one s & one p orbital, is different. It has the characteristics of both the mixing orbitals, and thus has an oval shape. The bigger lobes of sp² hybrid orbitals which result from the combination of one s & two p orbitals are pear shaped because of the greater contribution of two dumb-bell shaped p orbitals. In the case of sp³ hybrid orbitals which result from the mixing of one s & three p orbitals, the contribution of three p orbitals predominates over the contribution of one s orbital.

RULES FOR CONSRUCTING WAVE
FUNCTIONS FOR HYBRID ORBITALS
*1:- The Wave functions of hybrid orbitals are constructed by taking into consideration the linear combination of the wave function of the appropriate AOs. Thus, wave function of the ith hybrid orbital formed from s & p AOs is
Ψi = a i фs +bi фpx+ci фpy+di фpz
Where фs ,фpx ,фpy ,ф pz are the atomic orbital wave functions which constitute an orthogonal set .Two wave functions фm & фn
Are said to form a orthogonal set if
∫ фn фn dґ = 1
The coefficients of wave functions ai , bi ,ci ,di can be calculated by three conditions:-
• Each wave function of a hybrid orbital is normalized , that is,
∫ Ψi2 dґ = 1.It can be shown that for such a normalized wave function
ai2 +bi2 +ci2 +di2 = 1
• Each hybrid orbital in the set of hybrid orbitals is orthogonal to the other hybrid orbitals, i.e. ∫Ψ iΨjdΓ =0
From this relation, it can be shown that
aiaj +bibj +cicj +didj =0
• The squares of the coefficients of component wave-functions to the hybrid orbital wave functions over all the hybrid orbitals wherein they participate, equal unity, i.e.,
ai2 =1
This is because pure orbitals must be completely utilized in formation of hybrid orbitals.
*2:- Since the s orbital is spherically symmetric, each equivalent hybrid orbital of a set contains 1∕√n of the s orbital distributed in n orbitals , i.e. the coefficient of ф in each hybrid orbital is
1∕√n

DETERMINATION OF WAVE FUNCTIONS FOR THE sp HYBRID ORBITALS:-
Since one of the sp hybrid orbitals is to be placed along the +z axis, we shall invoke the hybridization of s, p ,p orbitals are in the XZ plane.The Wave functions for the three sp hybrid orbitals are:-

NUMERICALS ON 2ND ORDER REACTION

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 06:00:00 -0700

TOPIC: - NUMERICALS ON 2ND ORDER REACTION

NUMERICALS

Example 1: Butadiene dimerizes to form C8H12. This reaction is 2nd order in butadiene. If the rate constant for the reaction is 0.84 L/mol min, how long will it take for a 0.500 M sample of butadiene to dimerize until the butadiene concentration is 0.200 M?
Solution 1: Use the integrated rate law above
1/[C4H6] - 1/[C4H6]0 = kt
The intial concentration is 0.500 M, and the final concentration is 0.200 M. K is 0.84 L/mol min, so
1/0.200 mol/L - 1/0.500 mol/L = 0.84 L/mol min * t

t = 3.6 minutes

Example 2: A reaction 2A --> P has a second order rate law with k = 3.5E-4 L/mol-s. Calculate the time required for the concentration of A to change from 0.260 mol/L to 0.011 mol/L.

Solution 2: K1 =0.206 mol/l
K2 =0.011mol/l

1/ 0.011 - 1 / 0.260 = 3.5 x 10^-4 t

90.9 - 3.85 = 3.5 x 10^-4 t

t =248714 s = 69 h 5 min 14 s

Example 3: The rate constant for the SECOND ORDER reaction ; 2NO2 --> 2NO + O2 .Rate = 0.54 ? M.s at 300 Celsius. How long in seconds would it take for the concentration of NO2 to decrease from 0.62 M to 0.28 M?
Solution 3: time elapsed until the concentration has dropped from [NO2] to [NO2] is given by
t = ( 1/[NO2] - 1/[NO2] ) / k

Hence the time to decrease from 0.62M to 0.28M is
t = ( 1/.28M - 1/0.62M ) / 0.54M⁻¹s⁻¹
= 3.63S.
Example 4: Decomposition of a gas is a 2nd order reaction. It takes 40 min. for 40% of the gas to decompose. When its initial concentration is 0.04mol/l. Cal. Specific rate?
Solution 4: let the initial conc. Is a
 x=o.4a
 K=1/t * x/a(a-x)
 K=1/40* 0.4/0.6*0.04
 K=0.416mol/l/sec.

Example 5: a 2nd order reaction in which a=b is 20% comp. in 40 sec. How long will it take to comp. 60% ?
Solution 5: we know k=1/t* x/a(a-x)

 x=20* a/100
 x=1/160a…….( I )
now for x=60%
 k=1/t* 0.6a/a(a-0.6a)
 k=1/t* 3/2a…….( II )
from ( I ) ( II )
1/160a=1/t*3/2a
 t=240 sec.

Example 6: A second order reaction where a=b is 25% in 300sec. how long will it take for the reaction for 50 % completion?
Solution 6: a=100, T=300sec, x=25
a-x=75
now K=1/t* x/a(a-x)
k=1/300* 25/100(75)
 k=1/90,000
 t=1/k* x/a(a-x)
 t=1/1/90,000 * 50/5000
 t=900 sec.

Example 7: A second order reaction in which 25% of reaction comp. in 60 sec. How long will it take to comp. 75%.

Solution 7: a=100, t=600sec, a-x= 75
K=1/t * x/a(a-x)
K=1/60* 25/100(75)
K=1/18000
Now t=1/k* x/a(a-x)
T=18000* 75/2500
T=540 sec.

Example 8: 60% of the 2nd order reaction was completed in 60 min. when was it half comp.?
Solution 8: a=100, x=60, a-x=40
K=1/t*x/a(a-x)
K=1/60* 60/100(40)
K=0.00024
Now t=1/k* x/a(a-x)
T1/2= 1/ka
T1/2=0.024min.

REFERENCE
This topic has been referred from the following links :
1. ^ IUPAC Gold Book definition of rate law. See also: According to IUPAC Compendium of Chemical Terminology.

2. ^ Kenneth A. Connors Chemical Kinetics, the study of reaction rates in solution, 1991, VCH Publishers. This book contains all the rate equations in this article and their derivation.

3. ^ For a worked out example see: Determination of the Rotational Barrier for Kinetically Stable Conformational Isomers via NMR

4. 2D TLC An Introductory Organic Chemistry Experiment Gregory T. Rushton, William G.

5. Burns, Judi M. Lavin, Yong S. Chong, Perry Pellechia, and Ken D. Shimizu J. Chem. Educ. 2007, 84, 1499. Abstract

6. ^ José A. Manso et al."A Kinetic Approach to the Alkylating Potential of Carcinogenic Lactones" Chem. Res. Toxicol. 2005, 18, (7) 1161-1166

1. ^ Ruthenium(VI)-Catalyzed Oxidation of Alcohols by Hexacyanoferrate(III): An Example of Mixed Order Mucientes, Antonio E,; de la Peña, María A. J. Chem. Educ. 2006 83 1643. Abstract
• Chemical kinetics, reaction rate, and order (needs flash player)

• The reaction of crystal violet with sodium hydroxide: a kinetic study.

• Reaction kinetics, examples of important rate laws (lecture with audio).

• Rates of Reaction
Retrieved from "http://en.wikipedia.org/wiki/Order_of_reaction"

BUSINESS ENVIRONMET: Bajaj

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 05:59:00 -0700

CONTENTS:-

1. INTODUCTION
2. SOCIAL ENVIRONMENT
3. LEGAL ENVIRONMENT
4. POLITICAL ENVIRONMENT
5. ECONOMIC ENVIRONMENT
6. TECHNOLOGICAL ENVIRONMENT
7. SEGMENTATION
8. REFERENCES

BAJAJ INDUSTRIES
INTRODUCTION
Bajaj group of Industries, Nagpur started operations in 1961 under the able guidance of Late Shri Gangabisanji Bajaj. For last over 30 years it is being headed by Shri Hargovind Bajaj, An able Industrialist of region & founder Chairman of Vidarbha Industries Association. The prominent companies of the group include:The Company made its entry in the manufacture of the Cotton Ginning Machine and the first model was tested successfully in the year 1963. Company has emerged as a major success story with a steady growth rate, under the able guidance of Late Shri Gangabisan Bajaj and now under Shri Hargovind Bajaj.

RELATED PRODUCT INDUSTRIES

BAJAJ STEEL INDUSTRIES LTD
Bajaj Steel machines, cotton boll openers, cotton baling presses, pre-cleaners, and lint cleaners. The company also manufactures auto feeders and various there oequipments fours industries. or variIn addition, it offers steel bullock carts with break and tyres for sugar cane transportations. The company also operates in Bangladesh, Kenya, Madagascar, Malaysia, Myanmar, Nigeria, Sri Lanka, and Uganda. Bajaj Steel Industries is based in Nagpur, India. Industries, Ltd. manufactures cotton ginning machines in India. Its products include double roller ginning.

BAJAJ AUTO INDUSTRIES LTD
Bajaj Auto Ltd. is the largest exporter of two and three wheelers. With Kawasaki Heavy Industries of Japan, Bajaj manufactures state-of-the-art range of two-wheelers. The brand, Pulsar is continually dominating the Indian motorcycle market in the premium segment. Its Discover DTSi is also a successful bike on Indian roads. Founder -Jamnalal Bajaj, Year of Establishment -1926, Industry Automotive - Two & Three Wheelers, Business Group -The Bajaj Group. Since 1986, there is a technical tie-up of Bajaj Auto Ltd. with Kawasaki Heavy Industries of Japan to manufacture state-of-art range of latest two-wheelers in India. The JV has already given the Indian market the KB series, 4S and 4S Champion, Boxer, the Caliber series, and Wind125.

BAJAJ SUGAR INDUSTRIES LTD
Bajaj Hindustan Ltd., a part of the 'Bajaj Group', is India's number One sugar and ethanol manufacturing company. In December 2005, the company took over The Pratappur Sugar & Industries Limited in District Deoria, East U.P. This Plant, in operation since 1903, had a crushing capacity of 3,200 tcd, which has been increased to 6,000 tcd in the sugar season 2006-07. There is a large culturable area around this Unit with high potential of increasing cane cultivation and scaling up capacities significantly. The Pratappur Sugar & Industries Limited has now been renamed Bajaj Hindustan Sugar and Industries Limited (BHSIL) and is a subsidiary of BHL. With most sugar Plants in the vicinity closed, expansion of capacities at BHSIL will benefit the local farmer community and improve the rural economy of the region. This acquisition provides BHL a strategic foothold in the sugar-deficient region of Eastern UP and reaffirms the consolidation taking place in the sugar industry today. BHSIL has embarked upon significant new expansions. While the capacity of its existing sugar plant has been enhanced from 3,200 tcd to 6,000 tcd, three new sugar units are being set up in virgin, cane-rich areas of East UP to enhance its crushing capacity to 40,000 tcd over the next one year. BHSIL is also setting up a distillery with the capacity to manufacture 160 kilo liter per day of industrial alcohol/ ethanol. The Bajaj Hindustan Group is proud to be associated with the farming community of Uttar Pradesh and to contribute towards the economic development of rural India.

ALLIANZ GENERAL INSURANCE CO BAJAJ
Bajaj Allianz General Insurance Co in , India - Indians leading insurance company offers you Health Insurance, Travel Insurance, Home Insurance, Motor Insurance.
BAJAJ ELECTRICAL LTD
Bajaj Electricals Limited is the Consumer Goods and Lighting Company, manufacturing and trading these goods using very strong distribution network and setup at 21 Head office at Mumbai for IT department. Branches all over India and Head office at Mumbai.

SLEPT ANALYSIS
S-SOCIAL
L-LEGAL
E-ECONOMIC
P-POLITICAL
T-TECHNOLOGICAL

SOCIAL ENVIRONMENT OF BAJAJJ INDUSTRIES
Avnish Bajaj, founding managing director of venture capital fund has always been critical of especially the social networking business models, in India. In an interview to The Hindu Business Line, Bajaj has again raked up the issue of relevance of the social networking websites in the country. “I think people are wasting their time on Web 2.0 in India.” That’s a harsh statement on a lot of folks who are trying to develop businesses out of sites like Yaari.com, MingleBox and Desimartini, a new kid on the block.

Bajaj added: “People talk about the Internet being convenient, but it is not so in India. You need to go to a cyber café or you have to dial up a telephone line or use a slow broadband connection. Whereas in the US, 150 million households have broadband access all around the clock, sitting at home. When you have such a situation you can do social networking, but where is that happening in India? Do you think a person will go to a cyber café or any public environment to discuss everything about their life?” Bajaj’s point is that Indians may not need a platform for social networking yet. Indians culturally are not loners. He said: “Yes, there is a cultural barrier, as not many individuals will express themselves as in Myspace.com.” Besides, there are infrastructural barriers. So what is Bajaj’s point? “Fundamentally it is not about social networking but about community building. In India one needs to first create a product according to people’s needs and subsequently a community will form around it. An example would be Seventy mm, which solves a need. We are building a community product around it,” Bajaj concluded Bajaj in the community.Bajaj Auto is committed to nation-building and contributing to the uplift and development of the weaker sections of society. This is a legacy of our founders, Jamanalalji, Kamalnayanji and Ramkrishna Bajaj.
LEGAL BUSINESS ENVIRONMENT OF BAJAJ INDUSTRIES
Bajaj Auto Limited (BAL) is one of the leading players in the two-wheeler market, the second largest two-wheeler market in the world. The case traces the company's rise to dominance in the scooter segment of the market, and its eventual fall, against a backdrop of changes in customer tastes and preferences. It describes the reasons for the shift in demand and discusses the initiatives that the company undertook to regain lost ground. The case also discusses the competition in the Indian scooter market, and ends with a brief discussion on recent developments in the two-wheeler market Indian.
It was reported that the 125 cc motorbikes accounted for 55 percent of all bikes sold in India and the share would only grow in the future as the popularity of 100 cc bikes declined. While Bajaj planned to take on the entire market with its 125 cc motorbikes with DTSi technology, TVS was not far behind and was set to roll out as many as 12 products by the end of July 2008. Industry watchers felt that the legal tussle between the two arch rivals was not just about technology but the future of the two-wheeler market.
However, legal experts opined that it would not be easy for TVS to get Bajaj's patent revoked. Neither TVS nor any other competitor had objected to the granting of the patent to Bajaj for one full year though they had the right to do so. Now that the patent had been granted it would be very difficult for TVS to get the patent revoked, they said.

ECONOMICAL BUSINESS ENVIRONMENT OF BAJAJ INDUSTRIES
Bajaj Auto has got its strategy cut out. In a trying economic environment, it has plans to leverage its strong brands in the bigger (125cc +) motorcycle segment of the two-wheeler industry and the smaller three-wheeler segment of the light commercial vehicle industry.
This has worked for the country’s second largest two-wheeler manufacturer in the recent past. The advantages of this focus: Higher prices, a richer product mix, and lower cost of sales.
POLITICAL BUSINESS OF BAJAJ INDSTRIES
The world's largest two-wheeler manufacturer had reiterated its acknowledged market leadership by recording an impressive 22.4% growth in sales for the month of September 2008.

Going forward, with new launches on the anvil and this festive season, the company is confident of performing well in the coming months.
a winning product in each segment and increasing our touch points from 3,000 at present to 3,500 by the end of this year. The focus is on building strong brand identities through extensive campaigns and on ground marketing activities. It is also exploring previously untouched rural markets.
TECHNOLOGICAL ENVIRONMENT

COUPLE of farmers with Gandhi topics pedal their way leisurely over a bridge across the gurgling waters of the Indrayani river in Pune district. "We've just passed Sant Tukaram's birthplace," informs the driver, referring to the 16th century savant who hailed from this region. The late afternoon scene is almost pastoral as the car races through this piece of rural Maharashtra heading to Chakan where Bajaj Auto's latest two-wheeler plant has sprung up. Rocky terrain interspersed with lush, green fields fall behind as the plant nears. There on a rocky outcrop, Bajaj built its two-wheeler plant, which, it claims, is the most modern one in the industry.
Fountains spouting in a large pool of water and pleasing landscaping greet the visitor at the plant, which is built on two levels, taking advantage of the natural gradient.
Kicked off in year 2000, Bajaj's Chakan plant has a capacity to make six lakh two-wheelers a year, though only approximately 60 per cent of the capacity is being utilised now. The plant, which makes Bajaj's blockbuster bike, the Pulsar, will hum in full when Bajaj kicks off production in a few months of what it hopes will be a potential winner, a bike codenamed K60. Around 750 workers, all diploma-holders, you're told — 100 of them from Coimbatore in Tamil Nadu — produce the Pulsar and small numbers of its Spirit. The shop floors are air-cooled to not just provide a comfortable working environment but also to keep dust out from precision parts. After all, Bajaj wants to keep Pulsar at the top of the heap, considering that it has a 70 per cent market share of the premium end of the motorcycles market.
A

SEGMENTATION OF BAJAJ INDUSTRIES
Bajaj Auto launched its latest offering today. It is the most powerful sub-250cc bike on the Indian roads. Mr. Rajiv Bajaj, President, Bajaj Auto Ltd handed over the key to the first customer in the city in the premium bike segment Pulsar - the street sports bike in Mumbai.
Pulsar is all set to attract the performance seekers who today accounts for 5% of the total bike volume. But with a clever combination of three variants, Bajaj Pulsar has potential to sell 200,000 units in the next 24 months.
Targeted at the youth segment, the Bajaj Pulsar has been designed and styled as a mean masculine robust machine with dazzling looks and technically advanced mechanism that offers great performance. The symmetries between the muscular fuel tank, side panels and the rear panels give a very distinctive feel to Pulsar.

Business Profile

Bajaj Hindusthan Sugar & Industries produces sugar and ethanol. Formerly known as Pratappur Sugar & Industries, it was taken over by the Bajaj group in Dec. 2005 and renamed. Capacity of the plant has been increased to 6,000 TCD from existing 3200 TCD. The company is presently a subsidiary of Bajaj Hindustan.

Financials
For the last year ended Mar.31, 2006, the company`s net sales rose 32.87% from the previous year to Rs 647.2 million. Net loss increased 10.81 times to Rs 68.5 million. For the latest quarter ended Jun. 30, 2006, sales rose 26.91% to Rs 209.40 million compared with the corresponding quarter last year. The company also earned a profit of Rs 1.50 million as against a loss of Rs 20.20 million in last year's quarter.

REFERENCE:-
www.bajajngp.com/company.html
www.bajajindia.net/
www.surfindia.com/automobile/bajaj-auto-ltd.html
www.bajajcarpet.com/
wikimapia.org/416044/Bajaj-
business.mapsofindia.com/india-company/bajaj

AMPLIFIER

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 03:42:00 -0700

TERM PAPER
ON
AMPLIFIER

Submitted to: Submitted By:
Mam Suman Rani Sachin Dadhwal
Dept of Physics Roll No-R256A15
Regg No-10801248

Contents
1.Introduction
2.Working
3.Summary
4.Review of litrature
5.Bibleography

1.Introduction to Amplifiers
In "Electronics", signal amplifiers are widely used devices as they have the ability to amplify a relatively small input voltage signal, for example from a Sensor or microphone, into a much larger output signal to drive a Relay, lamp or loudspeaker for example. There are many forms of amplifiers, from Operational Amplifiers and Small Signal Amplifiers up to Large Signal and Power Amplifiers. Amplifiers can be thought of as a simple box or block containing the amplifying device, such as a Transistor, Field Effect Transistor or Op-amp, and which has two input terminals and two output terminals with the output signal being greater than that of the input signal, being "Amplified".
An amplifier has three main properties, Input Resistance or (Rin), Output Resistance or (Rout) and of course Gain or (A). No matter how complicated an amplifier circuit is, a general amplifier model can be used to show the relationship of these three properties.
Ideal Amplifier Model

Amplifier Gain
Then the gain of an amplifier can be said to be the relationship that exists between the signal measured at the output with the signal measured at the input. There are three different kinds of Amplifier Gain, Voltage Gain (Av), Current Gain (Ai) and Power Gain (Ap) and examples of these are given below.
Amplifier Gain of the Input Signal

Voltage Amplifier Gain

Current Amplifier Gain

Power Amplifier Gain

2. Working of amplifier

When people refer to "amplifiers," they're usually talking about stereo components or musical equipment. But this is only a small representation of the spectrum of audio amplifiers. There are actually amplifiers all around us. You'll find them in televisions, computers, portable CD players and most other devices that use a speaker to produce sound.
Sound is a fascinating phenomenon. When something vibrates in the atmosphere, it moves the air particles around it. Those air particles in turn move the air particles around them, carrying the pulse of the vibration through the air. Our ears pick up these fluctuations in air pressure and translate them into electrical signals the brain can process.
Electronic sound equipment works the same basic way. It represents sound as a varying electric current. Broadly speaking, there are three steps in this sort of sound reproduction:
Sound waves move a microphone diaphragm back and forth, and the microphone translates this movement into an electrical signal. The electrical signal fluctuates to represent the compressions and rarefactions of the sound wave.
A recorder encodes this electrical signal as a pattern in some sort of medium -- as magnetic impulses on tape, for example, or as grooves in a record.
A player (such as a tape deck) re-interprets this pattern as an electrical signal and uses this electricity to move a speaker cone back and forth. This re-creates the air-pressure fluctuations originally recorded by the microphone.
As you can see, all the major components in this system are essentially translators: They take the signal in one form and put it into another. In the end, the sound signal is translated back into its original form, a physical sound wave.
The theory of the laser amplifier is developed for conditions in which the strength of the input signal is increased from small values, where the amplification is linear, to larger values, where the amplification becomes nonlinear. The below-threshold laser amplifier oscillates at a single frequency equal to that of the input signal, and its properties are found by solution of the nonlinear equation of motion for the single excitation amplitude. For the above-threshold laser amplifier, the effects of the nonlinear behavior are to shift the laser frequency from its free-running value and to transfer intensity from the laser line to the signal frequency and to a range of satellite lines, whose frequency detunings are integer multiples of the signal detuning. The intensities of the various emission lines of the laser are calculated by power-series expansions of the field amplitudes up to terms of fourth order in the input signal strength. The onset of injection locking is determined by the conditions for which the intensity at the shifted free-running laser frequency falls to zero. The injection-locked state is characterized by a single excitation frequency equal to that of the input signal, and its properties are found by solution of the same nonlinear equation of motion as for the below-threshold amplifier. The ranges of input signal strength and detuning are determined for which the injection-locked state is stable. The energy conservation properties of the laser amplifier are considered for each of its operating states.
The theory of small-signal laser amplification is developed for homogeneously broadened systems in which the three main decay rates, for the collective atomic dipole moment, the population inversion, and the field in the laser cavity, have arbitrary relative magnitudes. The calculations extend previous work on class-A and -B lasers, where the dipole decay rate greatly exceeds the other two rates, to class-C lasers, where the dipole decay rate is comparable to the others. The free-running laser is assumed to excite a single longitudinal mode of the cavity, whose frequency generally differs from that of the coupled atomic transition. The linear gains of the laser are derived for input signals whose frequencies are close to that of the lasing mode or to one of its adjacent nonlasing longitudinal modes. Divergences in the gain that occur for these two arrangements are associated with the previously studied single-mode and multimode instabilities of the free-running laser, respectively.
3.Summary:

1.Resistor choice

The resistors used in these configurations are typically in the kΩ range. Resistors less than 1 kΩ cause excessive current flow and possible damage to the device. Resistors greater than 1 MΩ cause excessive thermal noise and make the circuit operation susceptible to significant errors due to otherwise negligibly small bias or leakage currents. Resistors adjacent to gaps between metal contacts will combine with parasitic capacitances to change the frequency-dependent circuit behavior. In general, while resistor choice may seem arbitrary below, in a real implementation, the size and placement of resistors affects the influence of small unmodeled quantities on the output.

2.Power supply effects

Although the power supplies are not shown in the operational amplifier designs below, they can be critical in operational amplifier design. Not only can they introduce pernicious effects, but they also can be useful in designs (e.g., the extend the power capability of an operational amplifier to effectively building a second output stage that delivers additional current based on the current demanded by the operational amplifier).

Operational amplifiers have a specified PSRR that indicates how well the output can reject signals that appear on the power supply inputs. Power supply inputs are often noisy in large designs because the power supply is used by nearly every component in the design, and inductance effects prevent current from being instantaneously delivered to every component at once. As a consequence, when a component requires large injections of current (e.g., a digital component that is frequently switching from one state to another), nearby components can experience sagging at their connection to the power supply. This problem can be mitigated with copious use of bypass capacitors placed connected across each power supply pin and ground. When bursts of current are required by a component, the component can bypass the power supply by receiving the current directly from the nearby capacitor (which is then slowly charged by the power supply).
3.Voltage follower:

Voltage followerUsed as a buffer amplifier, to eliminate loading effects or to interface impedances (connecting a device with a high source impedance to a device with a low input impedance). Due to the strong feedback, this circuit tends to get unstable when driving a high capacity load. This can be avoided by connecting the load through a resistor.
4.Summing amplifier

Sums several (weighted) voltages

• When , and Rf independent

• When

• Output is inverted
5.Comparator

Compares two voltages and switches its output to indicate which voltage is larger.
•
(where Vs is the supply voltage and the opamp is powered by + Vs and − Vs.)
6.Zero level detector
Voltage divider reference
• Zener sets reference voltage
4.Revuew of literature
I. INTRODUCTION
The limiting factor in processing a wide instantaneous dynamic range signal is the dynamic range of the processing circuitry. Logarithmic amplification solves this problem by transforming a large input dynamic range that increases logarithmically to a small output dynamic range which increases linearly. Since the compression ratio is known, the input signal amplitude information is preserved.
There are several types of logarithmic amplifiers presently available. These include logarithmic video amplifiers, detector logarithmic video amplifiers, logarithmic IF amplifiers, and true logarithmic amplifiers. Each type has its own unique properties. The comparison in Table I shows the key differences in specifications.

II. LOGARITHMIC VIDEO AMPLIFIER
A Log Video Amplifier (LVA) is a nonlinear video amplifier having the transfer characteristics shown in Figure 1.
This Transfer Function is of the form of:
V out = A*log (V in) + Vo
A = slope of the log video amp (mV/dBV)
Vin = input video signal (volts)
Vo = zero signal offset voltage (volts)

The above devices may be further distinguished as AC or DC coupled, "pulse on pulse," and dual (or extended range) types. Within th epulse on pulse category, there can be unipolar and bipolar devices. Often there are "accessory" functions included with the device design and packaging. These extra electrical functions may include threshold circuits, BITE circuits, and AGC or "null loop" functions.

A. AC Coupled Designs (earlier design techniques)

Historically, the most popular devices are the AC coupled ones. "AC" coupling means simply that there are one or more capacitors in series with the signal path through the device and thus DC inputs (or detected CW inputs) cannot be processed. AC coupled devices are the simplest structures and are often adequate performers since individual radar pulses are usually only a few microseconds in duration and of relatively low duty cycle. Such amplifiers can be designed with more than one circuit approach. The design technique used is shown in Figure 2.

In operation, a detected radar pulse is the input to a serial chain of pulse amplifier stages. These amplifier stages are fast, low power feedback amplifiers. Each stage is designed to smoothly limit at a specific signal level so that, for input signals increasing over the dynamic range, the fourth stage will limit first, then the third, and so on.
While there are other approaches to the realization of a log transfer characteristic, this exclusive approach is highly advantageous because it does not rely on precise semiconductor junction behavior for either transfer accuracy or stability. The "segments" are constructed at relatively high signal voltage levels and gains such that the break points are all easily set with resistor ratios. The involved resistors and reference voltages are all very stable, and volume production of a precisely performing product is readily attainable.
B. DC Coupled Designs (AMC techniques)
DC coupled designs are used when the amplifier input and output frequency responses needs to extend down to DC to process DC input signals or high duty cycle pulsed inputs. DC coupling eliminates the pulse droop seen on AC coupled LVAs caused by large pulse widths and base line variation caused by high duty cycles.
A basic block diagram is shown in Figure 3. This LVA is composed of a cascade of video amplifiers with one or more differential pairs off of the output of each amplifier. The video amplifiers provide high temperature stable, extremely low noise, video gain. This amplified video signal is then the input to the differential pairs. The differential pairs act as transconductance amplifiers. The outputs of all differential pairs are summed together on a "sum" line then transformed into a voltage and video amplified.
DC coupling adds considerable complexity to the design due to the extremely high Dc gain present at low input signal levels (gains can be greater then 1000). Temperature stability becomes difficult, though not impossible, to obtain due to this high gain.

III. DETECTOR LOGARITHMIC VIDEO AMPLIFIER
A detector logarithmic video amplifier (DLVA) is an amplifier that has input to output characteristics as described in Figure 4. For an input signal varying over many decades, the resultant output is a linear change over a much narrower range. The ratio between the input dynamic range and the output dynamic range is the compression ratio.
A block diagram of a DLVA is represented in Figure 5. A DLVA is a detector followed by a logarithmic video amplifier. The detector can be either a Schottky barrier diode or a zero bias tunnel diode. A Schottky diode requires bias and complex temperature compensation circuitry. The tunnel diode, on the other hand, is used in a zero bias configuration (biasing a tunnel diode might improve TSS by 2 dB). The temperature stability of a tunnel diode is much better than a Schottky diode while keeping a reasonable conversion factor, though not as high as the Schottky diode. The type of detector used is determined by the design of the LVA input stage. The Schottky diode requires a high LVA input impedance, and the tunnel diode requires a very low input impedance in order to control the input VSWR.
In basic operation the signal of pulse or CW is detected. The resultant video signal is fed into the LVA and amplified logarithmically. The LVA must have the ability to compensate for the diode non-linearities at high input power levels where the detector is moving from square law operation to linear operation and then into saturation. Also, any temperature related offsets and conversion factor variations caused by detector must be compensated by the LVA.

There are a few important design rules that must be followed in order to achieve the best performance from DLVA design.
A. Dynamic Range
The dynamic range of the DLVA is limited by the dynamic range of the LVA, the TSS level, and the magnitude and linearity of the conversion factor of the detector diode. The lower limit of the dynamic range is determined by the TSS and conversion factor of the detector, and by the sensitivity of the LVA. A logarithmically related input/output transfer function will not occur until the detected video signal is within the lower limit of the LVA’s logging range.
The TSS of the detector is a function of diode type, bias, and RF and video bandwidths. For a given diode, the TSS can be calculated as TSS per 2 MHz + 101og (Bv/2). The diode is connected to the LVA such that the effective noise bandwidth improves the TSS level. This improvement can be calculated as follows:

The new TSS level is now the diode TSS level minus the improvement. The start logging level can be calculated from the TSS level, detector conversion factor, and LVA sensitivity.
The maximum input signal level that can be processed is a function of the detector saturation level, detector conversion factor, and the upper input signal limit to the LVA. The upper limit is reached when either
the detector enters the saturation region, or when the detector supplies a video input signal to the LVA, which causes no change in the video output signal.
B. FREQUENCY FLATNESS
The variation of the detector VSWR and conversion factor due to frequency will appear on the output of the DLVA as an error. The worst case error caused by the frequency variation of the detector occurs when the VSWR ripple and conversion factor variations are added to the LVA maximum deviation from a true logarithmic response. LVAs do not exhibit any errors due to Rf frequency variations; only the detector determines errors. By improving the input match of the detector, this variation can be minimized, but care must be taken not to degrade the detector sensitivity and pulse response at the same time.
C. TEMPERATURE EFFECTS
Most DLVAs are required to operate over large temperature ranges and, as such, any errors due to temperature effects must be minimized. Temperature variations will have an effect on the detector conversion factor, noise floor, and on the gain of the video amplifiers. DC offset caused by these variations will be present in the output of a DC coupled DLVA. The temperature effects on a DC coupled DLVA are the most complicated to overcome due to the very high linear video gain associated with signals close to the start logging area in the LVA. With this high gain, any movement in the detector, or any DC offset shift in the first linear video amplifier due to temperature change, will be translated as a huge DC offset in the log output. Attention must be given in the design to achieve a large degree of temperature stability. The use of tunnel diode detectors is an easy way of achieving good detector temperature stability.
D. LINEARITY
Linearity of a DLVA can be explained as the deviation of the actual transfer function of a given DLVA from theoretical value of the logarithmic ratio. A noticeable linearity problem is the point at which the detector is moving from square law operation to linear and then to compression. This effect usually starts at the level of -20 to -10 dBm. To compensate for this effect, an additional logging stage is added in parallel at this point to increase the output onto the sum line. This effectively adds more gain to the upper logging stages in order to compensate for te compression of the detector conversion factor. This type of compensation is good for about 15 to 20 dB of input level above the detector break point; anything beyond this level must bel solved in a more complicated detector design.
E. VSWR
The input VSWR of a DLVA is determined by the match of the detector diode. The VSWR of the detector is worse at the transition point from square law to linear region. One way of improving this is by using the detector as a current source. This is done by connecting te detector into a virtual ground on the first linear video amplifier in the LVA (tunnel diode). Another way is to build a matching network at the input of the detector. This type of solution will result in loss of sensitivity and more restricted frequency bandwidth of the input RF signal.
F. EXTENDED DYNAMIC RANGE
The lower end of the dynamic range of a DLVA is determined by the minimum detected signal at the detector. A common value for TSS of a Schottky detector is around -50 dBm per MHz. This value is then modified by the RF and video bandwidths. In some applications, where a wider dynamic range is required, the lowest input signal level will be below the minimum detected signal. Then Rf gain must be introduced
to the system. Because the dynamic range of the DLVA and the detector is limited, the system will have to be built out of a number of LVAs and detectors. In this configuration, the signal has two paths: one is through the RF amplifier and DLVA 1, and at the saturation point of DLVA 1 the signal is processed through DLVA 2. The summation of the two video signals will give an extended range DLVA.
Careful attention should be given to the temperature, frequency, and gain performance of the RF amplifier since that will effect the overall linearity performance. Another possible error is te level at which the output switches between DLVA 1 and DLVA 2. This error is controlled by the same manner the detector saturation error is handled.
The frequency response of a DLVA is solely determined by the detector that is capable of multi-octave high frequency responses.
Depending on the application, the DLVA can be AC or DC coupled. The DC coupled DLVA is required, operating up to and including CW or in a system with very high duty cycle pulses.
Recent advances in the design of RF amplifiers opened the way to a more realizable IF successive detection approach that has a much superior pulse performance and, at the same time, supplies a limited RF output. Though these advances are being made, the DLVA is still the most convenient approach for extra wide bandwidths (multi-octave), and for frequency bands above 10 GHz, it is practically the only economical realizable approach.

References
• Baumeister, R. F., & Leary, M. R. (1995). The need to belong: Desire for interpersonal attachments as a fundamental human motivation. Psychological Bulletin 117, 497-529.
• Severin, Werner J., Tankard, James W., Jr., (1979). Communication Theories: Origins, Methods, Uses. New York: Hastings House, ISBN 0801317037
• Witzany, G. (2006 ) Plant Communication from Biosemiotic Perspective. Plant Signaling & Behavior 1(4): 169-178.
• Witzany, G. (2007 ). Applied Biosemiotics: Fungal Communication. In: Witzany, G. (Ed.) Biosemiotics in Transdisciplinary Contexts. *Helsinki. Umweb, pp 295-301.
• Wark, McKenzie 1997 The Virtual Republic Allen and Unwin St Leonards pp 22-9
• ^ Verhoeven CJM, van Staveren A, Monna GLE, Kouwenhoven MHL and Yildiz E (2003). Structured electronic design: negative feedback amplifiers. Boston/Dordrecht: Kluwer Academic. pp.10. ISBN 1-4020-7590-1, http://worldcat.org/isbn/1-4020-7590-1.
• ^ Robert S. Symons (1998). "Tubes: Still vital after all these years". IEEE Spectrum 35 (4): 52–63. doi:10.1109/6.666962.
• ^ OTB - Below 535, A Historical Review of Continuous Wave Radio Frequency Power Generators
Bibleography:
www.google.com
www.wikipedia.com
www.mamma.com

Modulation

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 03:41:00 -0700

INDEX:-

1 MODULATION

2 DAIGRAM OF MODULATION

3 AMPLITUDE MODULATION

4 DAIGRAM OF AMPLITUDE MODULATION

5 INTER MODULATION

6 DAIGRAM OF INTER MODULATION

Modulation

A message signal usually spreads over a range of frequencies,called the signal band width.That is why message signals are also called base band signals-representing the band of frequencies of the original signal.
Suppose we wish to transmit an electric signal in the audio frequencies(AF)range(20 Hz to 20kHz) over a long distance.we can not do it,as such because of the following reasons:-

1. Size of antenna or aerial:-an antenna or aerial is needed both for transmission and reception.each anteena should have a size comparable to the wavelength of the signal,(atleast λ/4 in size),so that time variation of the signal is properly sensed by the antenna.

For an audio frequency signal of frequency ν=15k Hz,the wave length, λ=c/ ν =(3*10^8/15*10^3)=20000m.the length of the antenna= λ/4=20000/4=5000m.to set up an antenna of vertical height of 5000 m is practically impossible to construct and operate.
If transmission frequency were raised to 1M Hz ,then
λ=c/ ν=(3*10^8)/(10^6)=300m
the length of antenna would be (300/4)=75m which is reasonable.therefore there is an urgent need of converting the information into higf frequencies before transmission.
2. Effective power radiation by antenna:-theortical studies we reveal that power P radiated from a linear antenna of length l is
Pα1/l^2

As high powers are needed for good transmission,l should be small i.e. antenna length would be small ,for which wavelength λ should be small or frequency should be high.
Thus,this factor also points ou t to the need of using high frequency before transmission.
3. Mixing up of signals from different transmitters:-
Suppose many people are talking at the same time.we just cannot make out who is talking what.similarly,when many transmitters are transmitting baseband information signal simultaneously ,they get mixed up and there is no way to distinguish between them.the possible solution is ,communication at high frequency and allotting a band of frequencies to each message signal for it’s transmission s
All the reasons explained above o that theren is no mixing.this is what is being done for different radios and T.V.broadcast stations.
All the things explained above suggest that there is need for transmission at high frequencies.this is achieved by a process,called modulation ,wherein we superimpose the low audio frequency baseband message or information signals(called the modulating signal)on a high frequency wave(called,the carrier wave).the resultant wave is called the modulated wave,which nis transmitted.
In the process of modulation,some specific characteristics of the carrier wave is varied in accordance with the information or message signal .the carrier wave may be
(1)continuous(sinusoidal)wave,or
(2)pulse which is discontinuous.

Modulation Diagram:-

F1

F2

F1 + F2

Summation of F1 and F2 (Linear)
In a Linear System, when one sinusoid is Superimposed upon another, neither sinusoid is affected, and no frequencies are generated.

F1 x F2

Modulation/Demodulation is a Nonlinear Process where Two Sinusoids are Multiplied (F1 x F2).
The Product of this Multiplication--in the Time Domain--is a wave whose Amplitude is:
a(X) = a(F1) X a(F2).
However, in the Frequency Domain there is an Addition and Subtraction of Frequencies, i.e.,
F1 + F2, F1 - F2, F2 + F1, F2 - F1, etc.
In practical terms, the Amplitude of one of the two sinusoids is held to a Constant Value; therefore, the useful results of this process are only the Derived Frequencies.
________________________________________
Modulation per se is used to impress a message (voice, image, data, etc.) on to a carrier wave for transmission. A bandlimited range of frequencies that comprise the message (baseband) is translated to a higher range of frequencies. The bandlimited message is preserved, i.e., every frequency in that message is scaled by a constant value

Amplitude-Modulation

In the design of an AM transmitter there are two ways to go:
(1) The low level generation of AM (DSB + Carrier) and the progressive amplification of that RF signal with the final stage being a Linear RF amplifier--Class AB.
In the case of the low level modulation approach, one could use either a 2 quadrant or 4 quadrant multiplier as the modulator.
TOP

The difference being: with the 2 quadrant multiplier, negative modulation of greater than 100% causes severe distortion as well as interference on adjacent bands. This is due to the carrier being cut off when the 2 quadrant multiplier cannot furnish any output for negative values of the modulating signals, hence the RF output becomes a pulsed spectrum.

If, however, a 4 quadrant multiplier is used, negative modulation starts to appear as a double sideband suppressed carrier--or in this case, reduced carrier.
(2) The second method is the progressive amplification of only the Carrier Wave with the output stage being, a more efficient, Class C (non-linear) RF amplifier; the modulation is introduced as a very high level audio signal at the final stage --more precisely, the positive plate supply of the RF "Final" Amplifier is made to vary as the modulation audio input signal.
The High Level modulation cannot handle negative modulation of greater than 100%. As with the 2 quadrant multiplier in the first approach, the carrier is cut off during negative peaks that exceed 100% negative modulation.
Most commercial AM and FM transmitter output stages--called "Finals"--use Class "C" amplifiers.
Other transmitters, like Television (visual), SSB, etc., use "Linear Amplifiers," Class AB1 or AB2, which are a combination of Class A and Class B (both being much less efficient than the Class C amplifier).

INTERMODULATION

Intermodulation is a Special Case where two (or more) sinusoids effect one another to produce undesired products, i.e., Unwanted Frequencies. Again, this can only occur when both waves share the same NonLinear device.

To Clarify: What is a Nonlinear Device? It is Any Active Device
[1]. In normal designs, radio receivers, Stereos, etc., Intermodulation is not a problem. However, when these systems are subjected to Excessive Signal Level Input the active devices in the "Front End" are driven out of their Linear Operating Regions--into or near--Saturation and/or Cutoff, where they become, in effect, "Modulators."
[1] Active Devices: Transistors, Diodes, ICs, etc.
Passive devices: Resistors, Capacitors, Inductors, etc.

Intermodulation
F1

F2

F1 + F2

Summation of F1 and F2 (Linear)

F2 x F1

F2 + F1

F2 - F1

Products Resulting from F1 and F2 (Nonlinear)

ELECTROLYSIS

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 03:39:00 -0700

Electrolysis
In chemistry and manufacturing, electrolysis is a method of separating chemically bonded elements and compounds by passing an electric current through them.
History
• 1800 - William Nicholson and Johann Ritter decomposed water into hydrogen and oxygen.
1807 - Potassium Permangenate was discovered by Sir Humphry Davy
Overviews
Electrolysis involves the passage of an electric current through, in general, an ionic substance that is either molten or dissolved in a suitable solvent, resulting in chemical reactions at the electrodes. The positive electrode is called the cathode, and the negative electrode is the anode.[1] To be useful for electrolysis, the electrodes need to be able to conduct electricity, and metal electrodes are generally used. Graphite electrodes and semiconductor electrodes are also used. An ionic compound, or a compound that reacts with the solvent to produce ions (such as an acid) is dissolved in an appropriate solvent, or an ionic compound is melted by heat. Then some free ions exist in the liquid. An electrical potential is applied between a pair of electrodes immersed in the liquid. Each electrode attracts ions that are of the opposite charge. Therefore, positively-charged ions (called cations) move towards the electron-emitting (negative) cathode, whereas negatively-charged ions (termed anions) move toward the positive anode. The energy required to separate the ions, and cause them to gather at the respective electrodes, is provided by an electrical power supply. At the electrodes, electrons are absorbed or released by the ions, forming a collection of the desired element or compound.
Oxidation of ions or neutral molecules can take place at the anode, and the reduction of ions or neutral molecules at the cathode. For example, it is possible to oxidize ferrous ions to ferric ions at the anode:
.
It is also possible to reduce ferricyanide ions to ferrocyanide ions at the cathode:
Neutral molecules can also react at either electrode. For example: p-Benzoquinone can be reduced to hydroquinone at the cathode:

In the last example, H + ions (hydrogen ions) also take part in the reaction, and are provided by an acid in the solution, or the solvent itself (water, methanol etc). Electrolysis reactions involving H + ions are fairly common in acidic solutions. In alkaline solutions, reactions involving OH − (hydroxide ions) are common.
The substances oxidised or reduced can also be the solvent (usually water) or the electrodes. It is possible to have electrolysis involving gases. For instance, fuel cells often use oxygen and hydrogen gases as reactants.

A demonstration electrochemical cell setup resembling the Daniell cell
The amount of electrical energy that must be added equals the change in Gibbs free energy of the reaction plus the losses in the system. The losses can (in theory) be arbitrarily close to zero, so the maximum thermodynamic efficiency equals the enthalpy change divided by the free energy change of the reaction. In most cases, the electric input is larger than the enthalpy change of the reaction, so some energy is released in the form of heat. In some cases, for instance, in the electrolysis of steam into hydrogen and oxygen at high temperature, the opposite is true. Heat is absorbed from the surroundings, and the heating value of the produced hydrogen is higher than the electric input.
The following technologies are related to electrolysis:
• Electrochemical cells, including the hydrogen fuel cell, use the reverse of this process.
• Gel electrophoresis is an electrolysis wherein the solvent is a gel: It is used to separate substances, such as DNA strands, based on their electrical charge.
Electrolysis of water
One important use of electrolysis of water is to produce hydrogen.
2H2O(l) → 2H2(g) + O2(g)
This has been suggested as a way of shifting society toward using hydrogen as an energy carrier for powering electric motors and internal combustion engines.elctrolysis of water can be observed by passing direct current from a battery or other DC power supply through a cup of water (in practice a salt water solution increases the reaction intensity making it easier to observe). Using platinum electrodes, hydrogen gas will be seen to bubble up at the cathode, and oxygen will bubble at the anode. If other metals are used as the anode, there is a chance that the oxygen will react with the anode instead of being released as a gas, or that the anode will dissolve. For example, using iron electrodes in a sodium chloride solution electrolyte, iron oxides will be produced at the anode. With zinc electrodes in a sodium chloride electrolyte, the anode will dissolve, producing zinc ions (Zn2+) in the solution, and no oxygen will be formed. When producing large quantities of hydrogen, the use of reactive metal electrodes can significantly contaminate the electrolytic cell - which is why iron electrodes are not usually used for commercial electrolysis. Electrodes made of stainless steel can be used because they will not react with the oxygen.
The energy efficiency of water electrolysis varies widely. The efficiency is a measure of what fraction of electrical energy used is actually contained within the hydrogen. Some of the electrical energy is converted to heat, a useless by-product. Some reports quote efficiencies between 50% and 70%[1] This efficiency is based on the Lower Heating Value of Hydrogen. The Lower Heating Value of Hydrogen is total thermal energy released when hydrogen is combusted minus the latent heat of vaporisation of the water. This does not represent the total amount of energy within the hydrogen, hence the efficiency is lower than a more strict definition. Other reports quote the theoretical maximum efficiency of electrolysis as being between 80% and 94%.[2]. The theoretical maximum considers the total amount of energy absorbed by both the hydrogen and oxygen. These values refer only to the efficiency of converting electrical energy into hydrogen's chemical energy. The energy lost in generating the electricity is not included. For instance, when considering a power plant that converts the heat of nuclear reactions into hydrogen via electrolysis, the total efficiency is more likely to be between 25% and 40%.
NREL found that a kilogram of hydrogen (roughly equivalent to a gallon of gasoline) could be produced by wind powered electrolysis for between $5.55 in the near term and $2.27 in the long term.[2]
About four percent of hydrogen gas produced worldwide is created by electrolysis, and normally used onsite. Hydrogen is used for the creation of ammonia for fertilizer via the Haber process, and converting heavy petroleum sources to lighter fractions via hydrocracking.
Pioneers of batteries:
• Alessandro Volta
More recently, electrolysis of heavy water was performed by Fleischmann and Pons in their famous experiment, resulting in anomalous heat generation and the controversial claim of cold fusion.
Faraday's laws of electrolysis
First law of electrolysis
In 1832, Michael Faraday reported that the quantity of elements separated by passing an electrical current through a molten or dissolved salt is proportional to the quantity of electric charge passed through the circuit. This became the basis of the first law of electrolysis:
Second law of electrolysis
Faraday also discovered that the mass of the resulting separated elements is directly proportional to the atomic masses of the elements when an appropriate integral divisor is applied. This provided strong evidence that discrete particles of matter exist as parts of the atoms of elements.
Industrial uses
• Production of aluminium, lithium, sodium, potassium, magnesium
• Production of hydrogen for hydrogen cars and fuel cells; high-temperature electrolysis is also used for this
• Coulometric techniques can be used to determine the amount of matter transformed during electrolysis by measuring the amount of electricity required to perform the electrolysis
• Production of chlorine and sodium hydroxide
• Production of sodium chlorate and potassium chlorate
• Production of perfluorinated organic compounds such as trifluoroacetic acid
Electrolysis has many other uses:
• Electrometallurgy is the process of reduction of metals from metallic compounds to obtain the pure form of metal using electrolysis. For example, sodium hydroxide in its molten form is separated by electrolysis into sodium and oxygen, both of which have important chemical uses. (Water is produced at the same time.)
• Anodization is an electrolytic process that makes the surface of metals resistant to corrosion. For example, ships are saved from being corroded by oxygen in the water by this process. The process is also used to decorate surfaces.
• A battery works by the reverse process to electrolysis. Humphry Davy found that lithium acts as an electrolyte and provides electrical energy.
• Production of oxygen for spacecraft and nuclear submarines.
• Electroplating is used in layering metals to fortify them. Electroplating is used in many industries for functional or decorative purposes, as in vehicle bodies and nickel coins.
• Production of hydrogen for fuel, using a cheap source of electrical energy.[3]
• Electrolytic Etching of metal surfaces like tools or knives with a permanent mark or logo.
Electrolysis is also used in the cleaning and preservation of old artifacts. Because the process separates the non-metallic particles from the metallic ones, it is very useful for cleaning old coins and even larger objects.

REFERENCES:-
www.google.com
www.wikipedia.com

CITIGROUP GLOBAL ORGANISATION

noreply@blogger.com (prasanjit) — Mon, 11 May 2009 03:38:00 -0700

CITIGROUP GLOBAL ORGANISATION

Citigroup Inc., operating as City (pronounced Siti), is a major American financial services company based in New York City, formed from the merger of Citicorp and Travelers Group on April 7, 1998. The company employs approximately 358,000 staff around the world, and holds over 200 million customer accounts in more than 100 countries. It is the world's largest bank by revenues as of 2008. It is a primary dealer in US Treasury securities and its stock has been a component of the Dow Jones Industrial Average since March 17, 1997. Its single largest shareholder is Prince Al-Waleed bin Talal of Saudi Arabia, who has a 4.9% stake.
As of December 11, 2007, Vikram Pandit is Citigroup's current CEO, while Sir Win Bischoff the current chairman.

Citigroup Inc.

Type
Public (NYSE: C, TYO: 8710)

Founded New York, New York, USA (1812)

Headquarters New York, New York, USA

Key people Sir Win Bischoff, Chairman
Vikram Pandit, CEO
Gary Crittenden, CFO

Industry
Financial services

Products
Consumer Banking
Corporate Banking
Investment Banking
Global Wealth Management
Investment Research
Private Equity

Market cap
▼ US$140 billion (2008)

Revenue
▲ US$159.2 billion (2007)
Net income
▼ US$3.62 billion (2007)
Total assets
▲ US$2.2 trillion (2007)
Total equity
▼ US$113.6 billion (2007)
Employees
358,000 (2007)
Website
http://www.citigroup.com/

Citigroup's world headquarters building,
399 Park Avenue New York city.

History:
Citigroup was formed on October 8, 1998 following the $140 billion merger of Citicorp and Travelers Group to create the world's largest financial services organization. The history of the company is, thus, divided into the history of several firms that over time amalgamated into Citicorp, a multinational banking corporation operating in more than 100 countries; or Travelers Group, whose businesses covered credit services, consumer finance, brokerage, and insurance. As such, the company history dates back to the founding of: the City Bank of New York (later Citibank) in 1812; Bank Handlowy in 1870; Smith Barney in 1873, Banamex in 1884; Salomon Brothers in 1910.

Citigroup New York City

Citicorp:
The history begins with the City Bank of New York, which was chartered by New York State on June 16, 1812, with $2 million of capital. Serving a group of New York merchants, the bank opened for business on September 14 of that year, and Samuel Osgood was elected as the first President of the company. The company's name was changed to The National City Bank of New York in 1865 after it joined the new U.S. national banking system, and it became the largest American bank by 1895. It became the first contributor to the Federal Reserve Bank of New York in 1913, and the following year it inaugurated the first overseas branch of a U.S. bank in Buenos Aires. The 1918 purchase of U.S. overseas bank International Banking Corporation helped it become the first American bank to surpass $1 billion in assets, and it became the largest commercial bank in the world in 1929. As it grew, the bank became a leading innovator in financial services, becoming the first major U.S. bank to offer compound interest on savings (1921); unsecured personal loans (1928); customer checking accounts (1936) and the negotiable certificate of deposit (1961).
The bank changed its name to The First National City Bank of New York in 1955, which was shortened to First National City Bank on the 150th anniversary of the company's foundation in 1962. The company organically entered the leasing and credit card sectors, and its introduction of USD.Certificates of deposit in London marked the first new negotiable instrument in market since 1888. Later to become MasterCard, the bank introduced its First National City Charge Service credit card - popularly known as the "Everything card"
Citicorp and Travelers merger:
On April 6, 1998, the merger between Citicorp and Travelers Group was announced to the world creating a $140 billion firm with assets of almost $700 billion. The deal would enable Travelers to market mutual funds and insurance to Citicorp's retail customers while giving the banking divisions access to an expanded client base of investors and insurance buyers.

The first logo (1998-2007) of the merged company, incorporating the Travelers' "red umbrella".

Although presented as a merger, the deal was actually more like a stock swap, with Travelers Group purchasing the entirety of Citicorp shares for $70 billion, and issuing 2.5 new Citigroup shares for each Citicorp share. Through this mechanism, existing shareholders of each company owned about half of the new firm. While the new company maintained Citicorp's "Citi" brand in its name, it adopted Travelers' distinctive "red umbrella" as the new corporate logo, which was used until 2007.
The chairmen of both parent companies, John Reed and Sandy Weill respectively, were announced as co-chairmen and co-CEOs of the new company, Citigroup, Inc., although the vast difference in management styles between the two immediately presented question marks over the wisdom of such a setup.
The remaining provisions of the Glass-Steagall Act - enacted following the Great Depression - forbade banks to merge with insurance underwriters, and meant Citigroup had between two and five years to divest any prohibited assets. However, Weill stated at the time of the merger that they believed "that over that time the legislation will change...we have had enough discussions to believe this will not be a problem". Indeed, the passing of the Gramm-Leach-Bliley Act in November 1999 vindicated Reed and Weill's views, opening the door to financial services conglomerates offering a mix of commercial banking, investment banking, insurance underwriting and brokerage.
Travelers spin off:

The current logo for Travelers Companies

The company spun off its Travelers Property and Casualty insurance underwriting business in 2002. The spin off was prompted by the insurance unit's drag on Citigroup stock price because Traveler's earnings were more seasonal and vulnerable to large disasters. It was also difficult to sell this kind of insurance directly to customers since most industrial customers are accustomed to purchasing insurance through a broker.
The Travelers Property Casualty Corporation merged with The St. Paul Companies Inc. in 2004 forming The St. Paul Travelers Companies. Citigroup retained the life insurance and annuities underwriting business; however, it sold those businesses to MetLife in 2005. Citigroup still heavily sells all forms of insurance, but it no longer underwrites insurance.
Despite their divesting Travelers Insurance, Citigroup retained Travelers' signature red umbrella logo as its own until February 2007, when Citigroup agreed to sell the logo back to St. Paul Travelers, which renamed itself Travelers Companies. Citigroup also decided to adopt the corporate brand "Citi" for itself and virtually all its subsidiaries, except Primerica and Banamex.
On April 11, 2007 Citigroup said it will eliminate 17,000 jobs, or about 5 percent of its workforce, in a broad restructuring designed to cut costs and bolster its long underperforming stock.[15]
On January 7, 2008 Citigroup announced that it is considering cutting 5 percent to 10 percent of its work force, which totals 327,000.

REAL ESTATE:

Citigroup Center, Chicago:
Citigroup's most famous office building is the Citigroup Center, a diagonal-roof skyscraper located in East Midtown, Manhattan, New York City, which despite popular belief is not the company's headquarters building. Citigroup has its headquarters across the street in an anonymous-looking building at 399 Park Avenue (the site of the original location of the City National Bank). The headquarters is outfitted with nine luxury dining rooms, with a team of private chefs preparing a different menu for each day. The management team is on the third and fourth floors above a Citibank branch. Citigroup also leases a building in the TriBeCa neighborhood in Manhattan at 388 Greenwich St, that serves as headquarters for its Investment and Corporate Banking operations and was the former headquarters of the Travelers Group.
Strategically, all of Citigroup's New York City real estate, excluding the company's Smith Barney division and Wall Street trading division, lies along the New York City Subway's IND Queens Boulevard Line, served by the E and V trains. Consequently, the company's Midtown buildings—including 787 Seventh Avenue, 666 Fifth Avenue, 399 Park Avenue, 485 Lexington, 153 East 53rd Street (Citigroup Center), and Citicorp Building in Long Island City, Queens, are all no more than two stops away from each other. In fact, every company building lies above or right across the street from a subway station served by the E or V.
Chicago also plays home to an architectural beauty operated by Citigroup. Citicorp Center has a series of curved archways at its peak, and sits across the street from major competitor ABN AMRO's ABN AMRO Plaza. It has a host of retail and dining facilities serving thousands of Metra customers daily via the Ogilvie Transportation Center.

Current news of citigroup global organisation:

September 23, 2008 6pm – 9pm Make Your Mark; Make a Difference - Ripple Leadership
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hosted by Women in Banking and Finance Citigroup Centre, Canada Square, London E14 5LB

Citi to divest BPO business - Citigroup Global Services:
The restructuring exercise at Citi may see its BPO business being transferred to a leading BPO in India. Though speculation is ripe yet we have very strong singals from people within the group that the deal may take shape as soon as 2-3 weeks.The deal size though unconfirmed by anyone, is put at a whopping $ 1 bn

Citi BPO sale put on hold :
The much speculated sale of Citigroup Global Services (CGSL) to Genpact has been put on indefinite hold. Since Vikram Pandit taking over as the CEO from Chuck Prince there have been many changes. This move is also attributed to the same. Some people who are familiar with Vikram’s style of functioning said that a hard look would be taken at the connect between global cost cutting (with job cuts) and all of its international outsourcing operations, before any decision can be arrived at. Sources also say that he delay will give Genpact time for a re- look at the $700 mn valuation which is believed to be on the higher side

Recruitment ad:

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Careers:
At Citigroup Global Services, we learn, grow, innovate, and enjoy – all while we work.
And we are looking for people with similar bent of mind. Think you share the same thoughts? Then this is the right place for you to kickstart your career.
Right from the word go in Citigroup Global Services, there’s no looking back. With every day, you will learn to create customer commitment, take initiative and create value to each function.
When you work with Citigroup Global Services, you will look forward to every working day.
At Citigroup Global Services, you will find a work culture that is professional yet informal. People address each other on a first name basis irrespective of hierarchy. Youngsters are encouraged to take a high degree of responsibility, even as they are rewarded for merit, innovation and initiative.

Our work culture is driven by our core values that put prime importance on Process Orientation and Culture; Focus on Quality and Productivity, Young and Well-Trained Workforce, a Friendly and Informal Atmosphere, Quality Focus and Innovation.
Human Capital
• Employee strength of 9000 with about 2000 being post-graduates (Masters/ CPA/ MBA)
• Strong “Employer of Choice” proposition
• “Equal opportunity employer”, with about 41% of its workforce being women
• Attractive source for talent for Citigroup worldwide; over 200 employees have joined various Citigroup entities globally
• Attrition managed to levels well below industry standards
• Innovation in HR
o Partnered with NMIMS and LIBA, premier business schools in Mumbai and Chennai, to conduct an annual, part-time MBA program exclusively for CGSL employees.
o Investment in Training: In house training academy CGSTA with over 1 million training hours per year.
o Structured career planning for employee
Apply now:
What's more, as you grow in a free boundary less environment where meritocracy prevails, you will learn to eliminate hierarchy. And lead through uncompromising compliance, integrity and professionalism. So what are you waiting for? Pump up the wheels and get
Citigroup Global Services is in one of the world’s most dynamic and growing business spaces. We bring you a range of career options in various business segments. Interested? Read about our current openings below.
Or , you can send your resumes to cgsl.careers@citigroup.com
S.No. Date Position Location
1 22 Sep Walk-in for Automation Testers on Saturday,27-Sept between 11-1pm
Mumbai, Mumbai Suburbs
2 16 Sep Walk-in-International Technical Voice Process on 22nd Sep Mon, 5pm-7pm
Mumbai
3 11 Sep Senior Manager USD Investigation
Mumbai, Mumbai
4 11 Sep Assistant Vice President Employee Relations (HR)
Mumbai, Mumbai Suburbs

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COMMENTS:
Citigroup is a major American financial services company based in New York City, form from the merger of Citicorp and traverlers corp having employs 358,000(app) staffing around the world.
It has covered credit services, consumer finance, brokerage and insurance.
As it was formed on Oct 8, 1998 following the $140 billion merger of citicorps and travelers group to create the world’s largest financial services organisation.it provides the excellent services to the customers.